--- base_model: stabilityai/stable-diffusion-xl-base-1.0 library_name: diffusers license: openrail++ tags: - text-to-image - text-to-image - diffusers-training - diffusers - lora - template:sd-lora - stable-diffusion-xl - stable-diffusion-xl-diffusers instance_prompt: a photo of sks dog widget: - text: A photo of sks dog in a bucket output: url: image_0.png - text: A photo of sks dog in a bucket output: url: image_1.png - text: A photo of sks dog in a bucket output: url: image_2.png - text: A photo of sks dog in a bucket output: url: image_3.png --- # SDXL LoRA DreamBooth - DKTech/dreambooth-test-1 ## Model description These are DKTech/dreambooth-test-1 LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0. The weights were trained using [DreamBooth](https://dreambooth.github.io/). LoRA for the text encoder was enabled: False. Special VAE used for training: madebyollin/sdxl-vae-fp16-fix. ## Trigger words You should use a photo of sks dog to trigger the image generation. ## Download model Weights for this model are available in Safetensors format. [Download](DKTech/dreambooth-test-1/tree/main) them in the Files & versions tab. ## Intended uses & limitations #### How to use Set up the environment on command-line / terminal. ```bash # Create and activate conda environment conda create –name dreambooth python=3.10 conda activate dreambooth # Install ipykernel (needed only if you want to run the inference inside a jupyter-notebook) conda install -c anaconda ipykernel python -m ipykernel install --user --name=dreambooth # Clone and install diffusers package git clone https://github.com/huggingface/diffusers cd diffusers pip install -e . # Browse to examples/dreambooth directory in the diffusers installation directory cd examples/dreambooth # Install dreambooth sdxl training dependencies pip install -r requirements_sdxl.txt ``` Run the inference in Python. ```python from huggingface_hub.repocard import RepoCard from diffusers import DiffusionPipeline import torch lora_model_id = "DKTech/dreambooth-test-1" card = RepoCard.load(lora_model_id) base_model_id = card.data.to_dict()["base_model"] pipe = DiffusionPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16) pipe = pipe.to("cuda") pipe.load_lora_weights(lora_model_id) image = pipe("A picture of an elephant that looks like a dog.", num_inference_steps=25).images[0] image.save("my_image.png") ``` #### Fine tuning the original model This model was created by fine tuning the original stable diffusion model based on the instructions here- https://github.com/huggingface/diffusers/blob/main/examples/dreambooth/README_sdxl.md Various other base models (other than stable diffusion) can also be fine tuned using DreamBooth. For example, some discussion on fine tuning Playground 2.5 model can be found here- https://github.com/huggingface/diffusers/pull/7126 #### Limitations and bias [TODO: provide examples of latent issues and potential remediations] ## Training details [TODO: describe the data used to train the model]