Update README.md
Browse files
README.md
CHANGED
|
@@ -1,25 +1,26 @@
|
|
| 1 |
---
|
| 2 |
language: en
|
| 3 |
tags:
|
| 4 |
-
-
|
| 5 |
-
-
|
| 6 |
-
|
|
|
|
| 7 |
datasets:
|
| 8 |
-
-
|
| 9 |
metrics:
|
| 10 |
-
-
|
| 11 |
-
---
|
| 12 |
|
| 13 |
# MaxCushion - SDXL Fine-tuned Model
|
| 14 |
|
| 15 |
-
This is a fine-tuned Stable Diffusion XL (SDXL) model based on [FLUX.1-dev](https://huggingface.co/black-forest-labs/FLUX.1-dev). It's designed to generate high-quality images with a focus on [specific theme or style your model specializes in]
|
| 16 |
|
| 17 |
## Model Details
|
| 18 |
|
| 19 |
- **Base Model:** [FLUX.1-dev](https://huggingface.co/black-forest-labs/FLUX.1-dev)
|
| 20 |
- **Type:** Stable Diffusion XL (SDXL)
|
| 21 |
- **Language(s):** English
|
| 22 |
-
- **License:** [
|
| 23 |
|
| 24 |
## Usage
|
| 25 |
|
|
@@ -30,37 +31,11 @@ from diffusers import StableDiffusionXLPipeline
|
|
| 30 |
import torch
|
| 31 |
|
| 32 |
model_id = "colt12/maxcushion"
|
| 33 |
-
pipe = StableDiffusionXLPipeline.from_pretrained(
|
|
|
|
|
|
|
| 34 |
pipe = pipe.to("cuda")
|
| 35 |
|
| 36 |
prompt = "Your prompt here"
|
| 37 |
image = pipe(prompt).images[0]
|
| 38 |
-
image.save("generated_image.png")
|
| 39 |
-
```
|
| 40 |
-
|
| 41 |
-
## Parameters
|
| 42 |
-
|
| 43 |
-
- `prompt`: The text prompt to generate an image from.
|
| 44 |
-
- `negative_prompt` (optional): Text prompt that the model should not use for image generation.
|
| 45 |
-
- `num_inference_steps` (default: 30): Number of denoising steps.
|
| 46 |
-
- `guidance_scale` (default: 7.5): How closely the model should follow the prompt.
|
| 47 |
-
|
| 48 |
-
## Examples
|
| 49 |
-
|
| 50 |
-
[Include a few example prompts and their resulting images here]
|
| 51 |
-
|
| 52 |
-
## Fine-tuning Details
|
| 53 |
-
|
| 54 |
-
This model was fine-tuned on [describe your dataset] using [describe your training process, e.g., tools, number of steps, learning rate, etc.].
|
| 55 |
-
|
| 56 |
-
## Limitations
|
| 57 |
-
|
| 58 |
-
[Describe any known limitations or biases of your model]
|
| 59 |
-
|
| 60 |
-
## Ethical Considerations
|
| 61 |
-
|
| 62 |
-
[Include any ethical considerations or guidelines for using your model]
|
| 63 |
-
|
| 64 |
-
## Contact
|
| 65 |
-
|
| 66 |
-
For questions or feedback, please [provide contact information or link to issues page].
|
|
|
|
| 1 |
---
|
| 2 |
language: en
|
| 3 |
tags:
|
| 4 |
+
- stable-diffusion-xl
|
| 5 |
+
- text-to-image
|
| 6 |
+
- diffusers
|
| 7 |
+
license: creativeml-openrail-m
|
| 8 |
datasets:
|
| 9 |
+
- your-dataset-name
|
| 10 |
metrics:
|
| 11 |
+
- your-metric
|
| 12 |
+
---
|
| 13 |
|
| 14 |
# MaxCushion - SDXL Fine-tuned Model
|
| 15 |
|
| 16 |
+
This is a fine-tuned Stable Diffusion XL (SDXL) model based on [FLUX.1-dev](https://huggingface.co/black-forest-labs/FLUX.1-dev). It's designed to generate high-quality images with a focus on **[specific theme or style your model specializes in]**.
|
| 17 |
|
| 18 |
## Model Details
|
| 19 |
|
| 20 |
- **Base Model:** [FLUX.1-dev](https://huggingface.co/black-forest-labs/FLUX.1-dev)
|
| 21 |
- **Type:** Stable Diffusion XL (SDXL)
|
| 22 |
- **Language(s):** English
|
| 23 |
+
- **License:** [CreativeML Open RAIL-M](https://huggingface.co/spaces/CompVis/stable-diffusion-license)
|
| 24 |
|
| 25 |
## Usage
|
| 26 |
|
|
|
|
| 31 |
import torch
|
| 32 |
|
| 33 |
model_id = "colt12/maxcushion"
|
| 34 |
+
pipe = StableDiffusionXLPipeline.from_pretrained(
|
| 35 |
+
model_id, torch_dtype=torch.float16, use_safetensors=True, variant="fp16"
|
| 36 |
+
)
|
| 37 |
pipe = pipe.to("cuda")
|
| 38 |
|
| 39 |
prompt = "Your prompt here"
|
| 40 |
image = pipe(prompt).images[0]
|
| 41 |
+
image.save("generated_image.png")
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|
|