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Opportunities and Challenges in High Pressure Processing of Foods
By Rastogi, N K; Raghavarao, K S M S; Balasubramaniam, V M; Niranjan, K; Knorr, D
Consumers increasingly demand convenience foods of the highest quality in terms of natural flavor and taste, and which are free from additives and preservatives. This demand has triggered the need for the development of a number of nonthermal approaches to food processing, of which high-pressure technology has proven to be very valuable. A number of recent publications have demonstrated novel and diverse uses of this technology. Its novel features, which include destruction of microorganisms at room temperature or lower, have made the technology commercially attractive. Enzymes and even spore forming bacteria can be inactivated by the application of pressure-thermal combinations, This review aims to identify the opportunities and challenges associated with this technology. In addition to discussing the effects of high pressure on food components, this review covers the combined effects of high pressure processing with: gamma irradiation, alternating current, ultrasound, and carbon dioxide or anti-microbial treatment. Further, the applications of this technology in various sectors-fruits and vegetables, dairy, and meat processing-have been dealt with extensively. The integration of high-pressure with other matured processing operations such as blanching, dehydration, osmotic dehydration, rehydration, frying, freezing / thawing and solid- liquid extraction has been shown to open up new processing options. The key challenges identified include: heat transfer problems and resulting non-uniformity in processing, obtaining reliable and reproducible data for process validation, lack of detailed knowledge about the interaction between high pressure, and a number of food constituents, packaging and statutory issues.
Keywords high pressure, food processing, non-thermal processing
Consumers demand high quality and convenient products with natural flavor and taste, and greatly appreciate the fresh appearance of minimally processed food. Besides, they look for safe and natural products without additives such as preservatives and humectants. In order to harmonize or blend all these demands without compromising the safety of the products, it is necessary to implement newer preservation technologies in the food industry. Although the fact that “high pressure kills microorganisms and preserves food” was discovered way back in 1899 and has been used with success in chemical, ceramic, carbon allotropy, steel/alloy, composite materials and plastic industries for decades, it was only in late 1980′s that its commercial benefits became available to the food processing industries. High pressure processing (HPP) is similar in concept to cold isostatic pressing of metals and ceramics, except that it demands much higher pressures, faster cycling, high capacity, and sanitation (Zimmerman and Bergman, 1993; Mertens and Deplace, 1993). Hite (1899) investigated the application of high pressure as a means of preserving milk, and later extended the study to preserve fruits and vegetables (Hite, Giddings, and Weakly, 1914). It then took almost eighty years for Japan to re- discover the application of high-pressure in food processing. The use of this technology has come about so quickly that it took only three years for two Japanese companies to launch products, which were processed using this technology. The ability of high pressure to inactivate microorganisms and spoilage catalyzing enzymes, whilst retaining other quality attributes, has encouraged Japanese and American food companies to introduce high pressure processed foods in the market (Mermelstein, 1997; Hendrickx, Ludikhuyze, Broeck, and Weemaes, 1998). The first high pressure processed foods were introduced to the Japanese market in 1990 by Meidi-ya, who have been marketing a line of jams, jellies, and sauces packaged and processed without application of heat (Thakur and Nelson, 1998). Other products include fruit preparations, fruit juices, rice cakes, and raw squid in Japan; fruit juices, especially apple and orange juice, in France and Portugal; and guacamole and oysters in the USA (Hugas, Garcia, and Monfort, 2002). In addition to food preservation, high- pressure treatment can result in food products acquiring novel structure and texture, and hence can be used to develop new products (Hayashi, 1990) or increase the functionality of certain ingredients. Depending on the operating parameters and the scale of operation, the cost of highpressure treatment is typically around US$ 0.05-0.5 per liter or kilogram, the lower value being comparable to the cost of thermal processing (Thakur and Nelson, 1998; Balasubramaniam, 2003).
The non-availability of suitable equipment encumbered early applications of high pressure. However, recent progress in equipment design has ensured worldwide recognition of the potential for such a technology in food processing (Could, 1995; Galazka and Ledward, 1995; Balci and Wilbey, 1999). Today, high-pressure technology is acknowledged to have the promise of producing a very wide range of products, whilst simultaneously showing potential for creating a new generation of value added foods. In general, high-pressure technology can supplement conventional thermal processing for reducing microbial load, or substitute the use of chemical preservatives (Rastogi, Subramanian, and Raghavarao, 1994).
Over the past two decades, this technology has attracted considerable research attention, mainly relating to: i) the extension of keeping quality (Cheftel, 1995; Farkas and Hoover, 2001), ii) changing the physical and functional properties of food systems (Cheftel, 1992), and iii) exploiting the anomalous phase transitions of water under extreme pressures, e.g. lowering of freezing point with increasing pressures (Kalichevsky, Knorr, and Lillford, 1995; Knorr, Schlueter, and Heinz, 1998). The key advantages of this technology can be summarized as follows:
1. it enables food processing at ambient temperature or even lower temperatures;
2. it enables instant transmittance of pressure throughout the system, irrespective of size and geometry, thereby making size reduction optional, which can be a great advantage;
3. it causes microbial death whilst virtually eliminating heat damage and the use of chemical preservatives/additives, thereby leading to improvements in the overall quality of foods; and
4. it can be used to create ingredients with novel functional properties.
The effect of high pressure on microorganisms and proteins/ enzymes was observed to be similar to that of high temperature. As mentioned above, high pressure processing enables transmittance of pressure rapidly and uniformly throughout the food. Consequently, the problems of spatial variations in preservation treatments associated with heat, microwave, or radiation penetration are not evident in pressure-processed products. The application of high pressure increases the temperature of the liquid component of the food by approximately 3C per 100 MPa. If the food contains a significant amount of fat, such as butter or cream, the temperature rise is greater (8-9C/100 MPa) (Rasanayagam, Balasubramaniam, Ting, Sizer, Bush, and Anderson, 2003). Foods cool down to their original temperature on decompression if no heat is lost to (or gained from) the walls of the pressure vessel during the holding stage. The temperature distribution during the pressure-holding period can change depending on heat transfer across the walls of the pressure vessel, which must be held at the desired temperature for achieving truly isothermal conditions. In the case of some proteins, a gel is formed when the rate of compression is slow, whereas a precipitate is formed when the rate is fast. High pressure can cause structural changes in structurally fragile foods containing entrapped air such as strawberries or lettuce. Cell deformation and cell damage can result in softening and cell serum loss. Compression may also shift the pH depending on the imposed pressure. Heremans (1995) indicated a lowering of pH in apple juice by 0.2 units per 100 MPa increase in pressure. In combined thermal and pressure treatment processes, Meyer (2000) proposed that the heat of compression could be used effectively, since the temperature of the product can be raised from 70-90C to 105-120C by a compression to 700 MPa, and brought back to the initial temperature by decompression.
As a thermodynamic parameter, pressure has far-reaching effects on the conformation of macromolecules, the transition temperature of lipids and water, and a number of chemical reactions (Cheftel, 1992; Tauscher, 1995). Phenomena that are accompanied by a decrease in volume are enhanced by pressure, and vice-versa (principle of Le Chatelier). Thus, under pressure, reaction equilibriums are shifted towards the most compact state, and the reaction rate constant is increased or decreased, depending on whether the “activation volume” of the reaction (i.e. volume of the activation complex less volume of reactants) is negative or positive. It is likely that pressure a\lso inhibits the availability of the activation energy required for some reactions, by affecting some other energy releasing enzymatic reactions (Farr, 1990). The compression energy of 1 litre of water at 400 MPa is 19.2 kJ, as compared to 20.9 kJ for heating 1 litre of water from 20 to 25C. The low energy levels involved in pressure processing may explain why covalent bonds of food constituents are usually less affected than weak interactions. Pressure can influence most biochemical reactions, since they often involve change in volume. High pressure controls certain enzymatic reactions. The effect of high pressure on protein/enzyme is reversible unlike temperature, in the range 100-400 MPa and is probably due to conformational changes and sub-unit dissociation and association process (Morild, 1981).
For both the pasteurization and sterilization processes, a combined treatment of high pressure and temperature are frequently considered to be most appropriate (Farr, 1990; Patterson, Quinn, Simpson, and Gilmour, 1995). Vegetative cells, including yeast and moulds, are pressure sensitive, i.e. they can be inactivated by pressures of ~300-600 MPa (Knorr, 1995; Patterson, Quinn, Simpson, and Gilmour, 1995). At high pressures, microbial death is considered to be due to permeabilization of cell membrane. For instance, it was observed that in the case of Saccharomyces cerevasia, at pressures of about 400 MPa, the structure and cytoplasmic organelles were grossly deformed and large quantities of intracellular material leaked out, while at 500 MPa, the nucleus could no longer be recognized, and a loss of intracellular material was almost complete (Farr, 1990). Changes that are induced in the cell morphology of the microorganisms are reversible at low pressures, but irreversible at higher pressures where microbial death occurs due to permeabilization of the cell membrane. An increase in process temperature above ambient temperature, and to a lesser extent, a decrease below ambient temperature, increases the inactivation rates of microorganisms during high pressure processing. Temperatures in the range 45 to 50C appear to increase the rate of inactivation of pathogens and spoilage microorganisms. Preservation of acid foods (pH ≤ 4.6) is, therefore, the most obvious application of HPP as such. Moreover, pasteurization can be performed even under chilled conditions for heat sensitive products. Low temperature processing can help to retain nutritional quality and functionality of raw materials treated and could allow maintenance of low temperature during post harvest treatment, processing, storage, transportation, and distribution periods of the life cycle of the food system (Knorr, 1995).
Bacterial spores are highly pressure resistant, since pressures exceeding 1200 MPa may be needed for their inactivation (Knorr, 1995). The initiation of germination or inhibition of germinated bacterial spores and inactivation of piezo-resistive microorganisms can be achieved in combination with moderate heating or other pretreatments such as ultrasound. Process temperature in the range 90-121C in conjunction with pressures of 500-800 MPa have been used to inactivate spores forming bacteria such as Clostridium botulinum. Thus, sterilization of low-acid foods (pH > 4.6), will most probably rely on a combination of high pressure and other forms of relatively mild treatments.
High-pressure application leads to the effective reduction of the activity of food quality related enzymes (oxidases), which ensures high quality and shelf stable products. Sometimes, food constituents offer piezo-resistance to enzymes. Further, high pressure affects only non-covalent bonds (hydrogen, ionic, and hydrophobic bonds), causes unfolding of protein chains, and has little effect on chemical constituents associated with desirable food qualities such as flavor, color, or nutritional content. Thus, in contrast to thermal processing, the application of high-pressure causes negligible impairment of nutritional values, taste, color flavor, or vitamin content (Hayashi, 1990). Small molecules such as amino acids, vitamins, and flavor compounds remain unaffected by high pressure, while the structure of the large molecules such as proteins, enzymes, polysaccharides, and nucleic acid may be altered (Balci and Wilbey, 1999).
High pressure reduces the rate of browning reaction (Maillard reaction). It consists of two reactions, condensation reaction of amino compounds with carbonyl compounds, and successive browning reactions including metanoidin formation and polymerization processes. The condensation reaction shows no acceleration by high pressure (5-50 MPa at 50C), because it suppresses the generation of stable free radicals derived from melanoidin, which are responsible for the browning reaction (Tamaoka, Itoh, and Hayashi, 1991). Gels induced by high pressure are found to be more glossy and transparent because of rearrangement of water molecules surrounding amino acid residues in a denatured state (Okamoto, Kawamura, and Hayashi, 1990).
The capability and limitations of HPP have been extensively reviewed (Thakur and Nelson, 1998; Smelt, 1998;Cheftal, 1995; Knorr, 1995; Fair, 1990; Tiwari, Jayas, and Holley, 1999; Cheftel, Levy, and Dumay, 2000; Messens, Van Camp, and Huyghebaert, 1997; Ontero and Sanz, 2000; Hugas, Garriga, and Monfort, 2002; Lakshmanan, Piggott,and Paterson, 2003; Balasubramaniam, 2003; Matser, Krebbers, Berg, and Bartels, 2004; Hogan, Kelly, and Sun, 2005; Mor-Mur and Yuste, 2005). Many of the early reviews primarily focused on the microbial efficacy of high-pressure processing. This review comprehensively covers the different types of products processed by highpressure technology alone or in combination with the other processes. It also discusses the effect of high pressure on food constituents such as enzymes and proteins. The applications of this technology in fruits and vegetable, dairy and animal product processing industries are covered. The effects of combining high- pressure treatment with other processing methods such as gamma- irradiation, alternating current, ultrasound, carbon dioxide, and anti microbial peptides have also been described. Special emphasis has been given to opportunities and challenges in high pressure processing of foods, which can potentially be explored and exploited.
EFFECT OF HIGH PRESSURE ON ENZYMES AND PROTEINS
Enzymes are a special class of proteins in which biological activity arises from active sites, brought together by a three- dimensional configuration of molecule. The changes in active site or protein denaturation can lead to loss of activity, or changes the functionality of the enzymes (Tsou, 1986). In addition to conformational changes, enzyme activity can be influenced by pressure-induced decompartmentalization (Butz, Koller, Tauscher, and Wolf, 1994; Gomes and Ledward, 1996). Pressure induced damage of membranes facilitates enzymesubstrate contact. The resulting reaction can either be accelerated or retarded by pressure (Butz, Koller, Tauscher, and Wolf, 1994; Gomes and Ledward, 1996; Morild, 1981). Hendrickx, Ludikhuy ze, Broeck, and Weemaes ( 1998) and Ludikhuyze, Van Loey, and Indrawati et al. (2003) reviewed the combined effect of pressure and temperature on enzymes related to the ity of fruits and vegetables, which comprises of kinetic information as well as process engineering aspects.
Pectin methylesterase (PME) is an enzyme, which normally tends to lower the viscosity of fruits products and adversely affect their texture. Hence, its inactivation is a prerequisite for the preservation of such products. Commercially, fruit products containing PME (e.g. orange juice and tomato products) are heat pasteurized to inactivate PME and prolong shelf life. However, heating can deteriorate the sensory and nutritional quality of the products. Basak and Ramaswamy (1996) showed that the inactivation of PME in orange juice was dependent on pressure level, pressure-hold time, pH, and total soluble solids. An instantaneous pressure kill was dependent only on pressure level and a secondary inactivation effect dependent on holding time at each pressure level. Nienaber and Shellhammer (2001) studied the kinetics of PME inactivation in orange juice over a range of pressures (400-600 MPa) and temperatures (25-5O0C) for various process holding times. PME inactivation followed a firstorder kinetic model, with a residual activity of pressure-resistant enzyme. Calculated D-values ranged from 4.6 to 117.5 min at 600 MPa/50C and 400 MPa/25C, respectively. Pressures in excess of 500 MPa resulted in sufficiently faster inactivation rates for economic viability of the process. Binh, Van Loey, Fachin, Verlent, Indrawati, and Hendrickx (2002a, 2002b) studied the kinetics of inactivation of strawberry PME. The combined effect of pressure and temperature on inactivation kinetics followed a fractional-conversion model. Purified strawberry PME was more stable toward high-pressure treatments than PME from oranges and bananas. Ly-Nguyen, Van Loey, Fachin, Verlent, Hendrickx (2002) showed that the inactivation of the banana PME enzyme during heating at temperature between 65 and 72.5C followed first order kinetics and the effect of pressure treatment of 600-700 MPa at 10C could be described using a fractionalconversion model. Stoforos, Crelier, Robert, and Taoukis (2002) demonstrated that under ambient pressure, tomato PME inactivation rates increased with temperature, and the highest rate was obtained at 75C. The inactivation rates were dramatically reduced as soon as the essing pressure was raised beyond 75C. High inactivation rates were obtained at a pressure higher than 700 MPa. Riahi and Ramaswamy (2003) studied high- pressure inactivation kinetics of PME isolated from a variety of sources and showed that PME from a microbial source was more resistant \to pressure inactivation than from orange peel. Almost a full decimal reduction in activity of commercial PME was achieved at 400 MPa within 20 min.
Verlent, Van Loey, Smout, Duvetter, Nguyen, and Hendrickx (2004) indicated that the optimal temperature for tomato pectinmethylesterase was shifted to higher values at elevated pressure compared to atmospheric pressure, creating the possibilities for rheology improvements by the application of high pressure.
Castro, Van Loey, Saraiva, Smout, and Hendrickx (2006) accurately described the inactivation of the labile fraction under mild-heat and high-pressure conditions by a fractional conversion model, while a biphasic model was used to estimate the inactivation rate constant of both the fractions at more drastic conditions of temperature/ pressure (10-64C, 0.1-800 MPa). At pressures lower than 300 MPa and temperatures higher than 54C, an antagonistic effect of pressure and temperature was observed.
Balogh, Smout, Binh, Van Loey, and Hendrickx (2004) observed the inactivation kinetics of carrot PME to follow first order kinetics over a range of pressure and temperature (650800 MPa, 10-40C). Enzyme stability under heat and pressure was reported to be lower in carrot juice and purified PME preparations than in carrots.
The presence of pectinesterase (PE) reduces the quality of citrus juices by destabilization of clouds. Generally, the inactivation of the enzyme is accomplished by heat, resulting in a loss of fresh fruit flavor in the juice. High pressure processing can be used to bypass the use of extreme heat for the processing of fruit juices. Goodner, Braddock, and Parish (1998) showed that the higher pressures (>600 MPa) caused instantaneous inactivation of the heat labile form of the enzyme but did not inactivate the heat stable form of PE in case of orange and grapefruit juices. PE activity was totally lost in orange juice, whereas complete inactivation was not possible in case of grapefruit juices. Orange juice pressurized at 700 MPa for l min had no cloud loss for more than 50 days. Broeck, Ludikhuyze, Van Loey, and Hendrickx (2000) studied the combined pressure-temperature inactivation of the labile fraction of orange PE over a range of pressure (0.1 to 900 MPa) and temperature (15 to 65C). Pressure and temperature dependence of the inactivation rate constants of the labile fraction was quantified using the well- known Eyring and Arrhenius relations. The stable fraction was inactivated at a temperature higher than 75C. Acidification (pH 3.7) enhanced the thermal inactivation of the stable fraction, whereas the addition of Ca^sup ++^ ions (IM) suppressed inactivation. At elevated pressure (up to 900 MPa), an antagonistic effect of pressure and temperature on inactivation of the stable fraction was observed. Ly-Nguyen, Van Loey, Smout, Ozean, Fachin, Verlent, Vu- Truong, Duvetter, and Hendrickx (2003) investigated the combined heat and pressure treatments on the inactivation of purified carrot PE, which followed a fractional-conversion model. The thermally stable fraction of the enzyme could not be inactivated. At a lower pressure (<300 MPa) and higher temperature (>50C), an antagonistic effect of pressure and heat was observed.
High pressures induced conformational changes in polygalacturonase (PG) causing reduced substrate binding affinity and enzyme inactivation. Eun, Seok, and Wan ( 1999) studied the effect of high-pressure treatment on PG from Chinese cabbage to prevent the softening and spoilage of plant-based foods such as kimchies without compromising quality. PG was inactivated by the application of pressure higher than 200 MPa for l min. Fachin, Van Loey, Indrawati, Ludikhuyze, and Hendrickx (2002) investigated the stability of tomato PG at different temperatures and pressures. The combined pressure temperature inactivation (300-600 MPa/50 -50C) of tomato PG was described by a fractional conversion model, which points to Ist-order inactivation kinetics of a pressure-sensitive enzyme fraction and to the occurrence of a pressure-stable PG fraction. Fachin, Smout, Verlent, Binh, Van Loey, and Hendrickx (2004) indicated that in the combination of pressure-temperature (5- 55C/100-600 MPa), the inactivation of the heat labile portion of purified tomato PG followed first order kinetics. The heat stable fraction of the enzyme showed pressure stability very similar to that of heat labile portion.
Peelers, Fachin, Smout, Van Loey, and Hendrickx (2004) demonstrated that effect of high-pressure was identical on heat stable and heat labile fractions of tomato PG. The isoenzyme of PG was detected in thermally treated (140C for 5 min) tomato pieces and tomato juice, whereas, no PG was found in pressure treated tomato juice or pieces.
Verlent, Van Loey, Smout, Duvetter, and Hendrickx (2004) investigated the effect of nigh pressure (0.1 and 500 MPa) and temperature (25-80C) on purified tomato PG. At atmospheric pressure, the optimum temperature for enzyme was found to be 55-60C and it decreased with an increase in pressure. The enzyme activity was reported to decrease with an increase in pressure at a constant temperature.
Shook, Shellhammer, and Schwartz (2001) studied the ability of high pressure to inactivate lipoxygenase, PE and PG in diced tomatoes. Processing conditions used were 400,600, and 800 MPa for 1, 3, and 5 min at 25 and 45C. The magnitude of the applied pressure had a significant effect in inactivating lipoxygenase and PG, with complete loss of activity occurring at 800 MPa. PE was very resistant to the pressure treatment.
Polyphenoloxidase and Pemxidase
Polyphenoloxidase (PPO) and peroxidase (POD), the enzymes responsible for color and flavor loss, can be selectively inactivated by a combined treatment of pressure and temperature. Gomes and Ledward (1996) studied the effects of pressure treatment (100-800 MPa for 1-20 min) on commercial PPO enzyme available from mushrooms, potatoes, and apples. Castellari, Matricardi, Arfelli, Rovere, and Amati ( 1997) demonstrated that there was a limited inactivation of grape PPO using pressures between 300 and 600 MPa. At 900 MPa, a low level of PPO activity was apparent. In order to reach complete inactivation, it may be necessary to use high- pressure processing treatments in conjunction with a mild thermal treatment (40-50C). Weemaes, Ludikhuyze, Broeck, and Hendrickx (1998) studied the pressure stabilities of PPO from apple, avocados, grapes, pears, and plums at pH 6-7. These PPO differed in pressure stability. Inactivation of PPO from apple, grape, avocado, and pear at room temperature (25C) became noticeable at approximately 600, 700, 800 and 900 MPa, respectively, and followed first-order kinetics. Plum PPO was not inactivated at room temperature by pressures up to 900 MPa. Rastogi, Eshtiaghi, and Knorr (1999) studied the inactivation effects of high hydrostatic pressure treatment (100-600 MPa) combined with heat treatment (0-60C) on POD and PPO enzyme, in order to develop high pressure-processed red grape juice having stable shelf-life. The studies showed that the lowest POD (55.75%) and PPO (41.86%) activities were found at 60C, with pressure at 600 and 100 MPa, respectively. MacDonald and Schaschke (2000) showed that for PPO, both temperature and pressure individually appeared to have similar effects, whereas the holding time was not significant. On the other hand, in case of POD, temperature as well as interaction between temperature and holding time had the greatest effect on activity. Namkyu, Seunghwan, and Kyung (2002) showed that mushroom PPO was highly pressure stable. Exposure to 600 MPa for 10 min reduced PPO activity by 7%; further exposure had no denaturing effect. Compression for 10 and 20 min up to 800 MPa, reduced activity by 28 and 43%, respectively.
Rapeanu, Van Loey, Smout, and Hendrickx (2005) indicated that the thermal and/or high-pressure inactivation of grape PPO followed first order kinetics. A third degree polynomial described the temperature/pressure dependence of the inactivation rate constants. Pressure and temperature were reported to act synergistically, except in the high temperature (≥45C)-low pressure (≥300 MPa) region where an antagonistic effect was observed.
Gomes, Sumner, and Ledward (1997) showed that the application of increasing pressures led to a gradual reduction in papain enzyme activity. A decrease in activity of 39% was observed when the enzyme solution was initially activated with phosphate buffer (pH 6.8) and subjected to 800 MPa at ambient temperature for 10 min, while 13% of the original activity remained when the enzyme solution was treated at 800 MPa at 60C for 10 min. In Tris buffer at pH 6.8 after treatment at 800 MPa and 20C, papain activity loss was approximately 24%. The inactivation of the enzyme is because of induced change at the active site causing loss of activity without major conformational changes. This loss of activity was due to oxidation of the thiolate ion present at the active site.
Weemaes, Cordt, Goossens, Ludikhuyze, Hendrickx, Heremans, and Tobback (1996) studied the effects of pressure and temperature on activity of 3 different alpha-amylases from Bacillus subtilis, Bacillus amyloliquefaciens, and Bacillus licheniformis. The changes in conformation of Bacillus licheniformis, Bacillus subtilis, and Bacillus amyloliquefaciens amylases occurred at pressures of 110, 75, and 65 MPa, respectively. Bacillus licheniformis amylase was more stable than amylases from Bacillus subtilis and Bacillus amyloliquefaciens to the combined heat/pressure treatment.
Riahi and Ramaswamy (2004) demonstrated that pressure inactivation of amylase in apple juice was significantly (P < 0.01 ) influenced by pH, pressure, holding time, and temperature. The inactivation was described using a bi-phasic model. The application of high pressure was sh\own to completely inactivate amylase. The importance of the pressure pulse and pressure hold approach for inactivation of amylase was also demonstrated.
High pressure denatures protein depending on the protein type, processing conditions, and the applied pressure. During the process of denaturation, the proteins may dissolve or precipitate on the application of high pressure. These changes are generally reversible in the pressure range 100-300 MPa and irreversible for the pressures higher than 300 MPa. Denaturation may be due to the destruction of hydrophobic and ion pair bonds, and unfolding of molecules. At higher pressure, oligomeric proteins tend to dissociate into subunits becoming vulnerable to proteolysis. Monomeric proteins do not show any changes in proteolysis with increase in pressure (Thakur and Nelson, 1998).
High-pressure effects on proteins are related to the rupture on non-covalent interactions within protein molecules, and to the subsequent reformation of intra and inter molecular bonds within or between the molecules. Different types of interactions contribute to the secondary, tertiary, and quaternary structure of proteins. The quaternary structure is mainly held by hydrophobic interactions that are very sensitive to pressure. Significant changes in the tertiary structure are observed beyond 200 MPa. However, a reversible unfolding of small proteins such as ribonuclease A occurs at higher pressures (400 to 800 MPa), showing that the volume and compressibility changes during denaturation are not completely dominated by the hydrophobic effect. Denaturation is a complex process involving intermediate forms leading to multiple denatured products. secondary structure changes take place at a very high pressure above 700 MPa, leading to irreversible denaturation (Balny and Masson, 1993).
Figure 1 General scheme for pressure-temperature phase diagram of proteins, (from Messens, Van Camp, and Huyghebaert, 1997).
When the pressure increases to about 100 MPa, the denaturation temperature of the protein increases, whereas at higher pressures, the temperature of denaturation usually decreases. This results in the elliptical phase diagram of native denatured protein shown in Fig. 1. A practical consequence is that under elevated pressures, proteins denature usually at room temperature than at higher temperatures. The phase diagram also specifies the pressure- temperature range in which the protein maintains its native structure. Zone III specifies that at high temperatures, a rise in denaturation temperature is found with increasing pressure. Zone II indicates that below the maximum transition temperature, protein denaturation occurs at the lower temperatures under higher pressures. Zone III shows that below the temperature corresponding to the maximum transition pressure, protein denaturation occurs at lower pressures using lower temperatures (Messens, Van Camp, and Huyghebaert, 1997).
The application of high pressure has been shown to destabilize casein micelles in reconstituted skim milk and the size distribution of spherical casein micelles decrease from 200 to 120 nm; maximum changes have been reported to occur between 150-400 MPa at 20C. The pressure treatment results in reduced turbidity and increased lightness, which leads to the formation of a virtually transparent skim milk (Shibauchi, Yamamoto, and Sagara, 1992; Derobry, Richard, and Hardy, 1994). The gels produced from high-pressure treated skim milk showed improved rigidity and gel breaking strength (Johnston, Austin, and Murphy, 1992). Garcia, Olano, Ramos, and Lopez (2000) showed that the pressure treatment at 25C considerably reduced the micelle size, while pressurization at higher temperature progressively increased the micelle dimensions. Anema, Lowe, and Stockmann (2005) indicated that a small decrease in the size of casein micelles was observed at 100 MPa, with slightly greater effects at higher temperatures or longer pressure treatments. At pressure >400 MPa, the casein micelles disintegrated. The effect was more rapid at higher temperatures although the final size was similar in all samples regardless of the pressure or temperature. At 200 MPa and 1O0C, the casein micelle size decreased slightly on heating, whereas, at higher temperatures, the size increased as a result of aggregation. Huppertz, Fox, and Kelly (2004a) showed that the size of casein micelles increased by 30% upon high-pressure treatment of milk at 250 MPa and micelle size dropped by 50% at 400 or 600 MPa.
Huppertz, Fox, and Kelly (2004b) demonstrated that the high- pressure treatment of milk at 100-600 MPa resulted in considerable solubilization of alphas 1- and beta-casein, which may be due to the solubilization of colloidal calcium phosphate and disruption of hydrophobic interactions. On storage of pressure, treated milk at 5C dissociation of casein was largely irreversible, but at 20C, considerable re-association of casein was observed. The hydration of the casein micelles increased on pressure treatment (100-600 MPa) due to induced interactions between caseins and whey proteins. Pressure treatment increased levels of alphas 1- and beta-casein in the soluble phase of milk and produced casein micelles with properties different to those in untreated milk. Huppertz, Fox, and Kelly (2004c) demonstrated that the casein micelle size was not influenced by pressures less than 200 MPa, but a pressure of 250 MPa increased the micelle size by 25%, while pressures of 300 MPa or greater, irreversibly reduced the size to 50% ofthat in untreated milk. Denaturation of alpha-lactalbumin did not occur at pressures less than or equal to 400 MPa, whereas beta-lactoglobulin was denatured at pressures greater than 100 MPa.
Galazka, Ledward, Sumner, and Dickinson (1997) reported loss of surface hydrophobicity due to application of 300 MPa in dilute solution. Pressurizing beta-lactoglobulin at 450 MPa for 15 minutes resulted in reduced solubility in water. High-pressure treatment induced extensive protein unfolding and aggregation when BSA was pressurized at 400 MPa. Beta-lactoglobulin appears to be more sensitive to pressure than alpha-lactalbumin. Olsen, Ipsen, Otte, and Skibsted (1999) monitored the state of aggregation and thermal gelation properties of pressure-treated beta-lactoglobulin immediately after depressurization and after storage for 24 h at 50C. A pressure of 150 MPa applied for 30 min, or pressures higher than 300 MPa applied for 0 or 30 min, led to formation of soluble aggregates. When continued for 30 min, a pressure of 450 MPa caused gelation of the 5% beta-lactoglobulin solution. Iametti, Tansidico, Bonomi, Vecchio, Pittia, Rovere, and DaIl’Aglio (1997) studied irreversible modifications in the tertiary structure, surface hydrophobicity, and association state of beta-lactoglobulin, when solutions of the protein at neutral pH and at different concentrations, were exposed to pressure. Only minor irreversible structural modifications were evident even for treatments as intense as 15 min at 900 MPa. The occurrence of irreversible modifications was time-dependent at 600 MPa but was complete within 2 min at 900 MPa. The irreversibly modified protein was soluble, but some covalent aggregates were formed. Subirade, Loupil, Allain, and Paquin (1998) showed the effect of dynamic high pressure on the secondary structure of betalactoglobulin. Thermal and pH sensitivity of pressure treated beta-lactoglobulin was different, suggesting that the two forms were stabilized by different electrostatic interactions. Walker, Farkas, Anderson, and Goddik (2004) used high- pressure processing (510 MPa for 10 min at 8 or 24C) to induce unfolding of beta-lactoglobulin and characterized the protein structure and surface-active properties. The secondary structure of the protein processed at 8C appeared to be unchanged, whereas at 24C alpha-helix structure was lost. Tertiary structures changed due to processing at either temperature. Model solutions containing the pressure-treated beta-lactoglobulin showed a significant decrease in surface tension. Izquierdo, Alli, Gmez, Ramaswamy, and Yaylayan (2005) demonstrated that under high-pressure treatments (100-300 MPa), the β-lactoglobulin AB was completely hydrolyzed by pronase and α-chymotrypsin. Hinrichs and Rademacher (2005) showed that the denaturation kinetics of beta-lactoglobulin followed second order kinetics while for alpha-lactalbumin it was 2.5. Alpha- lactalbumin was more resistant to denaturation than beta- lactoglobulin. The activation volume for denaturation of beta- lactoglobulin was reported to decrease with increasing temperature, and the activation energy increased with pressure up to 200 MPa, beyond which it decreased. This demonstrated the unfolding of the protein molecules.
Drake, Harison, Apslund, Barbosa-Canovas, and Swanson (1997) demonstrated that the percentage moisture and wet weight yield of cheese from pressure treated milk were higher than pasteurized or raw milk cheese. The microbial quality was comparable and some textural defects were reported due to the excess moisture content. Arias, Lopez, and Olano (2000) showed that high-pressure treatment at 200 MPa significantly reduced rennet coagulation times over control samples. Pressurization at 400 MPa led to coagulation times similar to those of control, except for milk treated at pH 7.0, with or without readjustment of pH to 6.7, which presented significantly longer coagulation times than their non-pressure treated counterparts.
Hinrichs and Rademacher (2004) demonstrated that the isobaric (200-800 MPa) and isothermal (-2 to 70C) denaturation of beta- lactoglobulin and alpha-lactalbumin of whey protein followed 3rd and 2nd order kinetics, respectively. Isothermal pressure denaturation of both beta-lactoglobulin A and B did not differ significantly and an increase in temperature resulted in an increase in thedenaturation rate. At pressures higher than 200 MPa, the denaturation rate was limited by the aggregation rate, while the pressure resulted in the unfolding of molecules. The kinetic parameters of denaturation were estimated using a single step non- linear regression method, which allowed a global fit of the entire data set. Huppertz, Fox, and Kelly (2004d) examined the high- pressure induced denaturation of alpha-lactalbumin and beta- lactoglobulin in dairy systems. The higher level of pressure- induced denaturation of both proteins in milk as compared to whey was due to the absence of casein micelles and colloidal calcium phosphate in the whey.
The conformation of BSA was reported to remain fairly stable at 400 MPa due to a high number of disulfide bonds which are known to stabilize its three dimensional structure (Hayakawa, Kajihara, Morikawa, Oda, and Fujio, 1992). Kieffer and Wieser (2004) indicated that the extension resistance and extensibility of wet gluten were markedly influenced by high pressure (up to 800 MPa), while the temperature and the duration of pressure treatment (30-80C for 2-20 min) had a relatively lesser effect. The application of high pressure resulted in a marked decrease in protein extractability due to the restructuring of disulfide bonds under high pressure leading to the incorporation of alpha- and gamma-gliadins in the glutenin aggregate. The change in secondary structure following high- pressure treatment was also reported.
The pressure treatment of myosin led to head-to-head interaction to form oligomers (clumps), which became more compact and larger in size during storage at constant pressure. Even after pressure treatment at 210 MPa for 5 minutes, monomieric myosin molecules increased and no gelation was observed for pressure treatment up to 210 MPa for 30 minutes. Pressure treatment did not also affect the original helical structure of the tail in the myosin monomers. Angsupanich, Edde, and Ledward (1999) showed that high pressure- induced denaturation of myosin led to formation of structures that contained hydrogen bonds and were additionally stabilized by disulphide bonds.
Application of 750 MPa for 20 minutes resulted in dimerization of metmyoglobin in the pH range of 6-10, whereas maximum pH was not at isoelectric pH (6.9). Under acidic pH conditions, no dimers were formed (Defaye and Ledward, 1995). Zipp and Kouzmann ( 1973) showed the formation of precipitate when pressurized (750 MPa for 20 minutes) near the isoelectric point, the precipitate redissolved slowly during storage. Pressure treatment had no effect on lipid oxidation in the case of minced meat packed in air at pressure less than 300 MPa, while the oxidation increased proportionally at higher pressures. However, on exposure to higher pressure, minced meat in contact with air oxidized rapidly. Pressures > 300-400 MPa caused marked denaturation of both myofibriller and sarcoplasmic proteins in washed pork muscle and pork mince (Ananth, Murano and Dickson, 1995). Chapleau and Lamballerie (2003) showed that high-pressure treatment induced a threefold increase in the surface hydrophobicity of myofibrillar proteins between O and 450 MPa. Chapleau, Mangavel, Compoint, and Lamballerie (2004) reported that high pressure modified the secondary structure of myofibrillar proteins extracted from cattle carcasses. Irreversible changes and aggregation were reported at a pressure higher than 300 MPa, which can potentially affect the functional properties of meat products. Lamballerie, Perron, Jung, and Cheret (2003) indicated that high pressure treatment increases cathepsin D activity, and that pressurized myofibrils are more susceptible to cathepsin D action than non- pressurized myofibrils. The highest cathepsin D activity was observed at 300 MPa. Cariez, Veciana, and Cheftel ( 1995) demonstrated that L color values increased significantly in meat treated at 200-350 MPa, the meat becoming pink, and a-value decreased in meat treated at 400-500 MPa to give a grey-brown color. The total extractable myoglobin decreased in meat treated at 200- 500 MPa, while the metmyoglobin content of meat increased and the oxymyoglobin decreased at 400500 MPa. Meat discoloration from pressure processing resulted in a whitening effect at 200-300 MPa due to globin denaturation, and/or haem displacement/release, or oxidation of ferrous myoglobin to ferric myoglobin at pressure higher than 400 MPa.
The conformation of the main protein component of egg white, ovalbumin, remains fairly stable when pressurized at 400 MPa, may be due to the four disulfide bonds and non-covalent interactions stabilizing the three dimensional structure of ovalbumin (Hayakawa, Kajihara, Morikawa, Oda, and Fujio, 1992). Hayashi, Kawamura, Nakasa and Okinada (1989) reported irreversible denaturation of egg albumin at 500-900 MPa with concomitant increase in susceptibility to subtilisin. Zhang, Li, and Tatsumi (2005) demonstrated that the pressure treatment (200-500 MPa) resulted in denaturation of ovalbumin. The surface hydrophobicity of ovalbumin was found to increase with increase in pressure treatment and the presence of polysaccharide protected the protein against denaturation. Iametti, Donnizzelli, Pittia, Rovere, Squarcina, and Bonomi (1999) showed that the addition of NaCl or sucrose to egg albumin prior to high- pressure treatment (up to 10 min at 800 MPa) prevented insolubulization or gel formation after pressure treatment. As a consequence of protein unfolding, the treated albumin had increased viscosity but retained its foaming and heat-gelling properties. Farr (1990) reported the modification of functionality of egg proteins. Egg yolk formed a gel when subjected to a pressure of 400 MPa for 30 minutes at 25C, kept its original color, and was soft and adhesive. The hardness of the pressure treated gel increased and adhesiveness decreased with an increase in pressure. Plancken, Van Loey, and Hendrickx (2005) showed that the application of high pressure (400- 700 MPa) to egg white solution resulted in an increase in turbidity, surface hydrophobicity, exposed sulfhydryl content, and susceptibility to enzymatic hydrolysis, while it resulted in a decrease in protein solubility, total sulfhydryl content, denaturation enthalpy, and trypsin inhibitory activity. The pressure- induced changes in these properties were shown to be dependent on the pressuretemperature and the pH of the solution. Speroni, Puppo, Chapleau, Lamballerie, Castellani, Aon, and Anton (2005) indicated that the application of high pressure (200-600 MPa) at 2OC to low- density lipoproteins did not change the solubility even if the pH is changed, whereas aggregation and protein denaturation were drastically enhanced at pH 8. Further, the application of high- pressure under alkaline pH conditions resulted in decreased droplet flocculation of low-density lipoproteins dispersions.
The minimum pressure required for the inducing gelation of soya proteins was reported to be 300 MPa for 10-30 minutes and the gels formed were softer with lower elastic modulus in comparison with heat-treated gels (Okamoto, Kawamura, and Hayashi, 1990). The treatment of soya milk at 500 MPa for 30 min changed it from a liquid state to a solid state, whereas at lower pressures and at 500 MPa for 10 minutes, the milk remained in a liquid state, but indicated improved emulsifying activity and stability (Kajiyama, Isobe, Uemura, and Noguchi, 1995). The hardness of tofu gels produced by high-pressure treatment at 300 MPa for 10 minutes was comparable to heat induced gels. Puppo, Chapleau, Speroni, Lamballerie, Michel, Anon, and Anton (2004) demonstrated that the application of high pressure (200-600 MPa) on soya protein isolate at pH 8.0 resulted in an increase in a protein hydorphobicity and aggregation, a reduction of free sulfhydryl content and a partial unfolding of the 7S and 11S fractions at pH 8. The change in the secondary structure leading to a more disordered structure was also reported. Whereas at pH 3.0, the protein was partially denatured and insoluble aggregates were formed, the major molecular unfolding resulted in decreased thermal stability, increased protein solubility, and hydorphobicity. Puppo, Speroni, Chapleau, Lamballerie, An, and Anton (2005) studied the effect of high- pressure (200, 400, and 600 MPa for 10 min at 10C) on the emulsifying properties of soybean protein isolates at pH 3 and 8 (e.g. oil droplet size, flocculation, interfacial protein concentration, and composition). The application of pressure higher than 200 MPa at pH 8 resulted in a smaller droplet size and an increase in the levels of depletion flocculation. However, a similar effect was not observed at pH 3. Due to the application of high pressure, bridging flocculation decreased and the percentage of adsorbed proteins increased irrespective of the pH conditions. Moreover, the ability of the protein to be adsorbed at the oil- water interface increased. Zhang, Li, Tatsumi, and Isobe (2005) showed that the application of high pressure treatment resulted in the formation of more hydrophobic regions in soy protein, which dissociated into subunits, which in some cases formed insoluble aggregates. High-pressure denaturation of beta-conglycinin (7S) and glycinin (11S) occurred at 300 and 400 MPa, respectively. The gels formed had the desirable strength and a cross-linked network microstructure.
Soybean whey is a by-product of tofu manufacture. It is a good source of peptides, proteins, oligosaccharides, and isoflavones, and can be used in special foods for the elderly persons, athletes, etc. Prestamo and Penas (2004) studied the antioxidative activity of soybean whey proteins and their pepsin and chymotrypsin hydrolysates. The chymotrypsin hydrolysate showed a higher antioxidative activity than the non-hydrolyzed protein, but the pepsin hydrolysate showed an opposite trend. High pressure processing at 100 MPa inc\reased the antioxidative activity of soy whey protein, but decreased the antioxidative activity of the hydrolysates. High pressure processing increased the pH of the protein hydrolysates. Penas, Prestamo, and Gomez (2004) demonstrated that the application of high pressure (100 and 200 MPa, 15 min, 37C) facilitated the hydrolysis of soya whey protein by pepsin, trypsin, and chymotrypsin. It was shown that the highest level of hydrolysis occurred at a treatment pressure of 100 MPa. After the hydrolysis, 5 peptides under 14 kDa with trypsin and chymotrypsin, and 11 peptides with pepsin were reported.
COMBINATION OF HIGHPRESSURE TREATMENT WITH OTHER NON-THERMAL PROCESSING METHODS
Many researchers have combined the use of high pressure with other non-thermal operations in order to explore the possibility of synergy between processes. Such attempts are reviewed in this section.
Crawford, Murano, Olson, and Shenoy (1996) studied the combined effect of high pressure and gamma-irradiation for inactivating Clostridium spmgenes spores in chicken breast. Application of high pressure reduced the radiation dose required to produce chicken meat with extended shelf life. The application of high pressure (600 MPa for 20 min at 8O0C) reduced the irradiation doses required for one log reduction of Clostridium spmgenes from 4.2 kGy to 2.0 kGy. Mainville, Montpetit, Durand, and Farnworth (2001) studied the combined effect of irradiation and high pressure on microflora and microorganisms of kefir. The irradiation treatment of kefir at 5 kGy and high-pressure treatment (400 MPa for 5 or 30 min) deactivated the bacteria and yeast in kefir, while leaving the proteins and lipids unchanged.
The exposure of microbial cells and spores to an alternating current (50 Hz) resulted in the release of intracellular materials causing loss or denaturation of cellular components responsible for the normal functioning of the cell. The lethal damage to the microorganisms enhanced when the organisms are exposed to an alternating current before and after the pressure treatment. High- pressure treatment at 300 MPa for 10 min for Escherichia coli cells and 400 MPa for 30 min for Bacillus subtalis spores, after the alternating current treatment, resulted in reduced surviving fractions of both the organisms. The combined effect was also shown to reduce the tolerant level of microorganisms to other challenges (Shimada and Shimahara, 1985, 1987; Shimada, 1992).
The pretreatment with ultrasonic waves (100 W/cm^sup 2^ for 25 min at 25C) followed by high pressure (400 MPa for 25 min at 15C) was shown to result in complete inactivation of Rhodoturola rubra. Neither ultrasonic nor high-pressure treatment alone was found to be effective (Knorr, 1995).
Carbon Dioxide and Argon
Heinz and Knorr (1995) reported a 3 log reduction of supercritical CO2 pretreated cultures. The effect of the pretreatment on germination of Bacillus subtilis endospores was monitored. The combination of high pressure and mild heat treatment was the most effective in reducing germination (95% reduction), but no spore inactivation was observed.
Park, Lee, and Park (2002) studied the combination of high- pressure carbon dioxide and high pressure as a nonthermal processing technique to enhance the safety and shelf life of carrot juice. The combined treatment of carbon dioxide (4.90 MPa) and high-pressure treatment (300 MPa) resulted in complete destruction of aerobes. The increase in high pressure to 600 MPa in the presence of carbon dioxide resulted in reduced activities of polyphenoloxidase (11.3%), lipoxygenase (8.8%), and pectin methylesterase (35.1%). Corwin and Shellhammer (2002) studied the combined effect of high-pressure treatment and CO2 on the inactivation of pectinmethylesterase, polyphenoloxidase, Lactobacillus plantarum, and Escherichia coli. An interaction was found between CO2 and pressure at 25 and 50C for pectinmethylesterase and polyphenoloxidase, respectively. The activity of polyphenoloxidase was decreased by CO2 at all pressure treatments. The interaction between CO2 and pressure was significant for Lactobacillus plantarum, with a significant decrease in survivors due to the addition of CO2 at all pressures studied. No significant effect on E. coli survivors was seen with CO2 addition. Truong, Boff, Min, and Shellhammer (2002) demonstrated that the addition of CO2 (0.18 MPa) during high pressure processing (600 MPa, 25C) of fresh orange juice increases the rate of PME inactivation in Valencia orange juice. The treatment time due to CO2 for achieving the equivalent reduction in PME activity was from 346 s to 111 s, but the overall degree of PME inactivation remained unaltered.
Fujii, Ohtani, Watanabe, Ohgoshi, Fujii, and Honma (2002) studied the high-pressure inactivation of Bacillus cereus spores in water containing argon. At the pressure of 600 MPa, the addition of argon reportedly accelerated the inactivation of spores at 20C, but had no effect on the inactivation at 40C.
The complex physicochemical environment of milk exerted a strong protective effect on Escherichia coli against high hydrostatic pressure inactivation, reducing inactivation from 7 logs at 400 MPa to only 3 logs at 700 MPa in 15 min at 20C. A substantial improvement in inactivation efficiency at ambient temperature was achieved by the application of consecutive, short pressure treatments interrupted by brief decompressions. The combined effect of high pressure (500 MPa) and natural antimicrobial peptides (lysozyme, 400 g/ml and nisin, 400 g/ml) resulted in increased lethality for Escherichia coli in milk (Garcia, Masschalck, and Michiels, 1999).
OPPORTUNITIES FOR HIGH PRESSURE ASSISTED PROCESSING
The inclusion of high-pressure treatment as a processing step within certain manufacturing flow sheets can lead to novel products as well as new process development opportunities. For instance, high pressure can precede a number of process operations such as blanching, dehydration, rehydration, frying, and solid-liquid extraction. Alternatively, processes such as gelation, freezing, and thawing, can be carried out under high pressure. This section reports on the use of high pressures in the context of selected processing operations.
Eshtiaghi and Knorr (1993) employed high pressure around ambient temperatures to develop a blanching process similar to hot water or steam blanching, but without thermal degradation; this also minimized problems associated with water disposal. The application of pressure (400 MPa, 15 min, 20C) to the potato sample not only caused blanching but also resulted in a four-log cycle reduction in microbial count whilst retaining 85% of ascorbic acid. Complete inactivation of polyphenoloxidase was achieved under the above conditions when 0.5% citric acid solution was used as the blanching medium. The addition of 1 % CaCl^sub 2^ solution to the medium also improved the texture and the density. The leaching of potassium from the high-pressure treated sample was comparable with a 3 min hot water blanching treatment (Eshtiaghi and Knorr, 1993). Thus, high- pressures can be used as a non-thermal blanching method.
Dehydration and Osmotic Dehydration
The application of high hydrostatic pressure affects cell wall structure, leaving the cell more permeable, which leads to significant changes in the tissue architecture (Fair, 1990; Dornenburg and Knorr, 1994, Rastogi, Subramanian, and Raghavarao, 1994; Rastogi and Niranjan, 1998; Rastogi, Raghavarao, and Niranjan, 2005). Eshtiaghi, Stute, and Knorr (1994) reported that the application of pressure (600 MPa, 15 min at 70C) resulted in no significant increase in the drying rate during fluidized bed drying of green beans and carrot. However, the drying rate significantly increased in the case of potato. This may be due to relatively limited permeabilization of carrot and beans cells as compared to potato. The effects of chemical pre-treatment (NaOH and HCl treatment) on the rates of dehydration of paprika were compared with products pre-treated by applying high pressure or high intensity electric field pulses (Fig. 2). High-pressure (400 MPa for 10 min at 25C) and high intensity electric field pulses (2.4 kV/cm, pulse width 300 s, 10 pulses, pulse frequency 1 Hz) were found to result in drying rates comparable with chemical pre-treatments. The latter pre-treatments, however, eliminated the use of chemicals (Ade- Omowaye, Rastogi, Angersbach, and Knorr, 2001).
Figure 2 (a) Effects of various pre-treatments such as hot water blanching, high pressure and high intensity electric field pulse treatment on dehydration characteristics of red paprika (b) comparison of drying time (from Ade-Omowaye, Rastogi, Angersbach, and Knorr, 2001).
Figure 3 (a) Variation of moisture and (b) solid content (based on initial dry matter content) with time during osmotic dehydration (from Rastogi and Niranjan, 1998).
Generally, osmotic dehydration is a slow process. Application of high pressures causes permeabilization of the cell structure (Dornenburg and Knorr, 1993; Eshtiaghi, Stute, and Knorr, 1994; Fair, 1990; Rastogi, Subramanian, and Raghavarao, 1994). This phenomenon has been exploited by Rastogi and Niranjan (1998) to enhance mass transfer rates during the osmotic dehydration of pineapple (Ananas comsus). High-pressure pre-treatments (100-800 MPa) were found to enhance both water removal as well as solid gain (Fig. 3). Measured diffusivity values for water were found to be four-fold greater, whilst solute (sugar) diffusivity values were found to be two-fold greater. Compression and decompression occurring during high pressure pre-treatment itself caused the removal of a significant amount of water, which was attributed to the cell wall rupture (Rastogi and Niranjan, 1998). Differential interference contrast microscopic examination showed the ext\ent of cell wall break-up with applied pressure (Fig. 4). Sopanangkul, Ledward, and Niranjan (2002) demonstrated that the application of high pressure (100 to 400 MPa) could be used to accelerate mass transfer during ingredient infusion into foods. Application of pressure opened up the tissue structure and facilitated diffusion. However, higher pressures above 400 MPa induced starch gelatinization also and hindered diffusion. The values of the diffusion coefficient were dependent on cell permeabilization and starch gelatinization. The maximum value of diffusion coefficient observed represented an eight-fold increase over the values at ambient pressure.
The synergistic effect of cell permeabilization due to high pressure and osmotic stress as the dehydration proceeds was demonstrated more clearly in the case of potato (Rastogi, Angersbach, and Knorr, 2000a, 2000b, 2003). The moisture content was reduced and the solid content increased in the case of samples treated at 400 MPa. The distribution of relative moisture (M/M^sub o^) and solid (S/S^sub o^) content as well as the cell permeabilization index (Zp) (shown in Fig. 5) indicate that the rate of change of moisture and solid content was very high at the interface and decreased towards the center (Rastogi, Angersbach, and Knorr, 2000a, 2000b, 2003).
Most dehydrated foods are rehydrated before consumption. Loss of solids during rehydration is a major problem associated with the use of dehydrated foods. Rastogi, Angersbach, Niranjan, and Knorr (2000c) have studied the transient variation of moisture and solid content during rehydration of dried pineapples, which were subjected to high pressure treatment prior to a two-stage drying process consisting of osmotic dehydration and finish-drying at 25C (Fig. 6). The diffusion coefficients for water infusion as well as for solute diffusion were found to be significantly lower in high-pressure pre- treated samples. The observed decrease in water diffusion coefficient was attributed to the permeabilization of cell membranes, which reduces the rehydration capacity (Rastogi and Niranjan, 1998). The solid infusion coefficient was also lower, and so was the release of the cellular components, which form a gel- network with divalent ions binding to de-esterified pectin (Basak and Ramaswamy, 1998; Eshtiaghi, Stute, and Knorr, 1994; Rastogi Angersbach, Niranjan, and Knorr, 2000c). Eshtiaghi, Stute, and Knorr (1994) reported that high-pressure treatment in conjunction with subsequent freezing could improve mass transfer during rehydration of dried plant products and enhance product quality.
Figure 4 Microstructures of control and pressure treated pineapple (a) control; (b) 300 MPa; (c) 700 MPa. ( 1 cm = 41.83 m) (from Rastogi and Niranjan, 1998).
Ahromrit, Ledward, and Niranjan (2006) explored the use of high pressures (up to 600 MPa) to accelerate water uptake kinetics during soaking of glutinous rice. The results showed that the length and the diameter the of the rice were positively correlated with soaking time, pressure and temperature. The water uptake kinetics was shown to follow the well-known Fickian model. The overall rates of water uptake and the equilibrium moisture content were found to increase with pressure and temperature.
Zhang, Ishida, and Isobe (2004) studied the effect of highpressure treatment (300-500 MPa for 0-380 min at 20C) on the water uptake of soybeans and resulting changes in their microstructure. The NMR analysis indicated that water mobility in high-pressure soaked soybean was more restricted and its distribution was much more uniform than in controls. The SEM analysis revealed that high pressure changed the microstructures of the seed coat and hilum, which improved water absorption and disrupted the individual spherical protein body structures. Additionally, the DSC and SDS-PAGE analysis revealed that proteins were partially denatured during the high pressure soaking. Ibarz, Gonzalez, Barbosa-Canovas (2004) developed the kinetic models for water absorption and cooking time of chickpeas with and without prior high-pressure treatment (275-690 MPa). Soaking was carried out at 25C for up to 23 h and cooking was achieved by immersion in boiling water until they became tender. As the soaking time increased, the cooking time decreased. High-pressure treatment for 5 min led to reductions in cooking times equivalent to those achieved by soaking for 60-90 min.
Ramaswamy, Balasubramaniam, and Sastry (2005) studied the effects of high pressure (33, 400 and 700 MPa for 3 min at 24 and 55C) and irradiation (2 and 5 kGy) pre-treatments on hydration behavior of navy beans by soaking the treated beans in water at 24 and 55C. Treating beans under moderate pressure (33 MPa) resulted in a high initial moisture uptake (0.59 to 1.02 kg/kg dry mass) and a reduced loss of soluble materials. The final moisture content after three hours of soaking was the highest in irradiated beans (5 kGy) followed by high-pressure treatment (33 MPa, 3 min at 55C). Within the experimental range of the study, Peleg’s model was found to satisfactorily describe the rate of water absorption of navy beans.
A reduction of 40% in oil uptake during frying was observed, when thermally blanched frozen potatoes were replaced by high pressure blanched frozen potatoes. This may be due to a reduction in moisture content caused by compression and decompression (Rastogi and Niranjan, 1998), as well as the prevalence of different oil mass transfer mechanisms (Knorr, 1999).
Solid Liquid Extraction
The application of high pressure leads to rearrangement in tissue architecture, which results in increased extractability even at ambient temperature. Extraction of caffeine from coffee using water could be increased by the application of high pressure as well as increase in temperature (Knorr, 1999). The effect of high pressure and temperature on caffeine extraction was compared to extraction at 100C as well as atmospheric pressure (Fig. 7). The caffeine yield was found to increase with temperature at a given pressure. The combination of very high pressures and lower temperatures could become a viable alternative to current industrial practice.
Figure 5 Distribution of (a, b) relative moisture and (c, d) solid content as well as (e, f) cell disi |
Notes on the Bible, by Albert Barnes, , at sacred-texts.com
Now the word of the Lord - , literally, "And, ..." This is the way in which the several inspired writers of the Old Testament mark that what it was given them to write was united onto those sacred books which God had given to others to write, and it formed with them one continuous whole. The word, "And," implies this. It would do so in any language, and it does so in Hebrew as much as in any other. As neither we, nor any other people, would, without any meaning, use the word, And, so neither did the Hebrews. It joins the four first books of Moses together; it carries on the history through Joshua, Judges, the Books of Samuel and of the Kings. After the captivity, Ezra and Nehemiah begin again where the histories before left off; the break of the captivity is bridged over; and Ezra, going back in mind to the history of God's people before the captivity, resumes the history, as if it had been of yesterday, "And in the first year of Cyrus." It joins in the story of the Book of Ruth before the captivity, and that of Esther afterward. At times, even prophets employ it, in using the narrative form of themselves, as Ezekiel, "and it was in the thirtieth year, in the fourth month, in the fifth day of the month, and I was in the captivity by the river of Chebar, the heavens opened and I saw." If a prophet or historian wishes to detach his prophecy or his history, he does so; as Ezra probably began the Book of Chronicles anew from Adam, or as Daniel makes his prophecy a whole by itself. But then it is the more obvious that a Hebrew prophet or historian, when he does begin with the word, "And," has an object in so beginning; he uses an universal word of all languages in its uniform meaning in all language, to join things together.
And yet more precisely; this form, "and the word of the Lord came to - saying," occurs over and over again, stringing together the pearls of great price of God's revelations, and uniting this new revelation to all those which had preceded it. The word, "And," then joins on histories with histories, revelations with revelations, uniting in one the histories of God's works and words, and blending the books of Holy Scripture into one divine book.
But the form of words must have suggested to the Jews another thought, which is part of our thankfulness and of our being Act 11:18, "then to the Gentiles also hath God given repentance unto life." The words are the self-same familiar words with which some fresh revelation of God's will to His people had so often been announced. Now they are prefixed to God's message to the pagan, and so as to join on that message to all the other messages to Israel. Would then God deal thenceforth with the pagan as with the Jews? Would they have their prophets? Would they be included in the one family of God? The mission of Jonah in itself was an earnest that they would, for God. Who does nothing fitfully or capriciously, in that He had begun, gave an earnest that He would carry on what He had begun. And so thereafter, the great prophets, Isaiah, Jeremiah, Ezekiel, were prophets to the nations also; Daniel was a prophet among them, to them as well as to their captives.
But the mission of Jonah might, so far, have been something exceptional. The enrolling his book, as an integral part of the Scriptures, joining on that prophecy to the other prophecies to Israel, was an earnest that they were to be parts of one system. But then it would be significant also, that the records of God's prophecies to the Jews, all embodied the accounts of their impenitence. Here is inserted among them an account of God's revelation to the pagan, and their repentance. "So many prophets had been sent, so many miracles performed, so often had captivity been foreannounced to them for the multitude of their sins. and they never repented. Not for the reign of one king did they cease from the worship of the calves; not one of the kings of the ten tribes departed from the sins of Jeroboam? Elijah, sent in the Word and Spirit of the Lord, had done many miracles, yet obtained no abandonment of the calves. His miracles effected this only, that the people knew that Baal was no god, and cried out, "the Lord He is the God." Elisha his disciple followed him, who asked for a double portion of the Spirit of Elijah, that he might work more miracles, to bring back the people.
He died, and, after his death as before it, the worship of the calves continued in Israel. The Lord marveled and was weary of Israel, knowing that if He sent to the pagan they would bear, as he saith to Ezekiel. To make trial of this, Jonah was chosen, of whom it is recorded in the Book of Kings that he prophesied the restoration of the border of Israel. When then he begins by saying, "And the word of the Lord came to Jonah," prefixing the word "And," he refers us back to those former things, in this meaning. The children have not hearkened to what the Lord commanded, sending to them by His servants the prophets, but have hardened their necks and given themselves up to do evil before the Lord and provoke Him to anger; "and" therefore "the word of the Lord came to Jonah, saying, Arise and go to Nineveh that great city, and preach unto her," that so Israel may be shewn, in comparison with the pagan, to be the more guilty, when the Ninevites should repent, the children of Israel persevered in unrepentance."
Jonah the son of Amittai - Both names occur here only in the Old Testament, Jonah signifies "Dove," Amittai, "the truth of God." Some of the names of the Hebrew prophets so suit in with their times, that they must either have been given them propheticly, or assumed by themselves, as a sort of watchword, analogous to the prophetic names, given to the sons of Hosea and Isaiah. Such were the names of Elijah and Elisha, "The Lord is my God," "my God is salvation." Such too seems to be that of Jonah. The "dove" is everywhere the symbol of "mourning love." The side of his character which Jonah records is that of his defect, his want of trust in God, and so his unloving zeal against those, who were to be the instruments of God against his people. His name perhaps preserves that character by which he willed to be known among his people, one who moaned or mourned over them.
Arise, go to Nineveh, that great city - The Assyrian history, as far as it has yet been discovered, is very bare of events in regard to this period. We have as yet the names of three kings only for 150 years. But Assyria, as far as we know its history, was in its meridian. Just before the time of Jonah, perhaps ending in it, were the victorious reigns of Shalmanubar and Shamasiva; after him was that of Ivalush or Pul, the first aggressor upon Israel. It is clear that this was a time Of Assyrian greatness: since God calls it "that great city," not in relation to its extent only, but its power. A large weak city would not have been called a "great city unto God" Jon 3:3.
And cry against it - The substance of that cry is recorded afterward, but God told to Jonah now, what message he was to cry aloud to it. For Jonah relates afterward, how he expostulated now with God, and that his expostulation was founded on this, that God was so merciful that He would not fulfill the judgment which He threatened. Faith was strong in Jonah, while, like Apostles "the sons of thunder," before the Day of Pentecost, he knew not" what spirit he was of." Zeal for the people and, as he doubtless thought, for the glory of God, narrowed love in him. He did not, like Moses, pray Exo 32:32, "or else blot me also out of Thy book," or like Paul, desire even to be "an anathema from Christ" Rom 9:3 for his people's sake, so that there might be more to love his Lord. His zeal was directed, like that of the rebuked Apostles, against others, and so it too was rebuked. But his faith was strong. He shrank back from the office, as believing, not as doubting, the might of God. He thought nothing of preaching, amid that multitude of wild warriors, the stern message of God. He was willing, alone, to confront the violence of a city of 600,000, whose characteristic was violence. He was ready, at God's bidding, to enter what Nahum speaks of as a den of lions Nah 2:11-12; "The dwelling of the lions and the feeding-place of the young lions, where the lion did tear in pieces enough for his whelps, and strangled for his lionesses." He feared not the fierceness of their lion-nature, but God's tenderness, and lest that tenderness should be the destruction of his own people.
Their wickedness is come up before Me - So God said to Cain, Gen 4:10. "The voice of thy brother's blood crieth unto Me from the ground:" and of Sodom Gen 18:20 :21, "The cry of Sodom and Gomorrah is great, because their sin is very grievous; the cry of it is come up unto Me." The "wickedness" is not the mere mass of human sin, of which it is said Jo1 5:19, "the whole world lieth in wickedness," but evil-doing toward others. This was the cause of the final sentence on Nineveh, with which Nahum closes his prophecy, "upon whom hath not thy wickedness passed continually?" It bad been assigned as the ground of the judgment on Israel through Nineveh Hos 10:14-15. "So shall Bethel do unto you, on account of the wickedness of your wickedness." It was the ground of the destruction by the flood Gen 6:5. "God saw that the wickedness of man was great upon the earth." God represents Himself, the Great Judge, as sitting on His Throne in heaven, Unseen but All-seeing, to whom the wickedness and oppressiveness of man against man "goes up," appealing for His sentence against the oppressor. The cause seems ofttimes long in pleading. God is long-suffering with the oppressor too, that if so be, he may repent. So would a greater good come to the oppressed also, if the wolf became a lamb. But meanwhile, " every iniquity has its own voice at the hidden judgment seat of God." Mercy itself calls for vengeance on the unmerciful.
But (And) Jonah rose up to flee ... from the presence of the Lord - literally "from being before the Lord." Jonah knew well, that man could not escape from the presence of God, whom he knew as the Self-existing One, He who alone is, the Maker of heaven, earth and sea. He did not "flee" then "from His presence," knowing well what David said Psa 139:7, Psa 139:9-10, "whither shall I go from Thy Spirit, or whither shall I flee from Thy presence? If I take the wings of the morning, and dwell in the uttermost parts of the sea, even there shall Thy hand lead me and Thy right hand shall hold me." Jonah fled, not from God's presence, but from standing before him, as His servant and minister. He refused God's service, because, as he himself tells God afterward Jon 4:2, he knew what it would end in, and he misliked it.
So he acted, as people often do, who dislike God's commands. He set about removing himself as far as possible from being under the influence of God, and from the place where he "could" fulfill them. God commanded him to go to Nineveh, which lay northeast from his home; and he instantly set himself to flee to the then furthermost west. Holy Scripture sets the rebellion before us in its full nakedness. "The word of the Lord came unto Jonah, go to Nineveh, and Jonah rose up;" he did something instantly, as the consequence of God's command. He "rose up," not as other prophets, to obey, but to disobey; and that, not slowly nor irresolutely, but "to flee, from" standing "before the Lord." He renounced his office. So when our Lord came in the flesh, those who found what He said to be "hard sayings," went away from Him, "and walked no more with Him" Joh 6:66. So the rich "young man went away sorrowful Mat 19:22, for he had great possessions."
They were perhaps afraid of trusting themselves in His presence; or they were ashamed of staying there, and not doing what He said. So men, when God secretly calls them to prayer, go and immerse themselves in business; when, in solitude, He says to their souls something which they do not like, they escape His Voice in a throng. If He calls them to make sacrifices for His poor, they order themselves a new dress or some fresh sumptuousness or self-indulgence; if to celibacy, they engage themselves to marry immediately; or, contrariwise, if He calls them not to do a thing, they do it at once, to make an end of their struggle and their obedience; to put obedience out of their power; to enter themselves on a course of disobedience. Jonah, then, in this part of his history, is the image of those who, when God calls them, disobey His call, and how He deals with them, when he does not abandon them. He lets them have their way for a time, encompasses them with difficulties, so that they shall "flee back from God displeased to God appeased."
"The whole wisdom, the whole bliss, the whole of man lies in this, to learn what God wills him to do, in what state of life, calling, duties, profession, employment, He wills him to serve Him." God sent each one of us into the world, to fulfill his own definite duties, and, through His grace, to attain to our own perfection in and through fulfilling them. He did not create us at random, to pass through the world, doing whatever self-will or our own pleasure leads us to, but to fulfill His will. This will of His, if we obey His earlier calls, and seek Him by prayer, in obedience, self-subdual, humility, thoughtfulness, He makes known to each by His own secret drawings, and, in absence of these, at times by His Providence or human means. And then , "to follow Him is a token of predestination." It is to place ourselves in that order of things, that pathway to our eternal mansion, for which God created us, and which God created for us.
So Jesus says Joh 10:27-28, "My sheep hear My voice and I know them, and they follow Me, and I give unto them eternal life, and they shall never perish, neither shall any man pluck them out of My Hand." In these ways, God has foreordained for us all the graces which we need; in these, we shall be free from all temptations which might be too hard for us, in which our own special weakness would be most exposed. Those ways, which people choose out of mere natural taste or fancy, are mostly those which expose them to the greatest peril of sin and damnation. For they choose them, just because such pursuits flatter most their own inclinations, and give scope to their natural strength and their moral weakness. So Jonah, disliking a duty, which God gave him to fulfill, separated himself from His service, forfeited his past calling, lost, as far as in him lay, his place among "the goodly fellowship of the prophets," and, but for God's overtaking grace, would have ended his days among the disobedient. As in Holy Scripture, David stands alone of saints, who had been after their calling, bloodstained; as the penitent robber stands alone converted in death; as Peter stands singly, recalled after denying his Lord; so Jonah stands, the one prophet, who, having obeyed and then rebelled, was constrained by the overpowering providence and love of God, to return and serve Him.
"Being a prophet, Jonah could not be ignorant of the mind of God, that, according to His great Wisdom and His unsearchable judgments and His untraceable and incomprehensible ways, He, through the threat, was providing for the Ninevites that they should not suffer the things threatened. To think that Jonah hoped to hide himself in the sea and elude by flight the great Eye of God, were altogether absurd and ignorant, which should not be believed, I say not of a prophet, but of no other sensible person who had any moderate knowledge of God and His supreme power. Jonah knew all this better than anyone, that, planning his flight, he changed his place, but did not flee God. For this could no man do, either by hiding himself in the bosom of the earth or depths of the sea or ascending (if possible) with wings into the air, or entering the lowest hell, or encircled with thick clouds, or taking any other counsel to secure his flight.
This, above all things and alone, can neither be escaped nor resisted, God. When He willeth to hold and grasp in His Hand, He overtaketh the swift, baffleth the intelligent, overthroweth the strong, boweth the lofty, tameth rashness, subdueth might. He who threatened to others the mighty Hand of God, was not himself ignorant of nor thought to flee, God. Let us not believe this. But since he saw the fall of Israel and perceived that the prophetic grace would pass over to the Gentiles, he withdrew himself from the office of preaching, and put off the command." "The prophet knoweth, the Holy Spirit teaching him, that the repentance of the Gentiles is the ruin of the Jews. A lover then of his country, he does not so much envy the deliverance of Nineveh, as will that his own country should not perish. - Seeing too that his fellow-prophets are sent to the lost sheep of the house of Israel, to excite the people to repentance, and that Balaam the soothsayer too prophesied of the salvation of Israel, he grieveth that he alone is chosen to be sent to the Assyrians, the enemies of Israel, and to that greatest city of the enemies where was idolatry and ignorance of God. Yet more he feared lest they, on occasion of his preaching, being converted to repentance, Israel should be wholly forsaken. For he knew by the same Spirit whereby the preaching to the Gentiles was entrusted to him, that the house of Israel would then perish; and he feared that what was at one time to be, should take place in his own time." "The flight of the prophet may also be referred to that of man in general who, despising the commands of God, departed from Him and gave himself to the world, where subsequently, through the storms of ill and the wreck of the whole world raging against him, he was compelled to feel the presence of God, and to return to Him whom he had fled. Whence we understand, that those things also which men think for their good, when against the will of God, are turned to destruction; and help not only does not benefit those to whom it is given, but those too who give it, are alike crushed. As we read that Egypt was conquered by the Assyrians, because it helped Israel against the will of God. The ship is emperiled which had received the emperiled; a tempest arises in a calm; nothing is secure, when God is against us."
Tarshish - , named after one of the sons of Javan, Gen 10:4. was an ancient merchant city of Spain, once proverbial for its wealth (Psa 72:10. Strabo iii. 2. 14), which supplied Judaea with silver Jer 10:9, Tyre with "all manner of riches," with iron also, tin, lead. Eze 27:12, Eze 27:25. It was known to the Greeks and Romans, as (with a harder pronunciation) Tartessus; but in our first century, it had either ceased to be, or was known under some other name. Ships destined for a voyage, at that time, so long, and built for carrying merchandise, were naturally among the largest then constructed. "Ships of Tarshish" corresponded to the "East-Indiamen" which some of us remember. The breaking of "ships of Tarshish by the East Wind" Psa 48:7 is, on account of their size and general safety, instanced as a special token of the interposition of God.
And went down to Joppa - Joppa, now Jaffa (Haifa), was the one well-known port of Israel on the Mediterranean. There the cedars were brought from Lebanon for both the first and second temple Ch2 3:16; Ezr 2:7. Simon the Maccabee (1 Macc. 14:5) "took it again for a haven, and made an entrance to the isles of the sea." It was subsequently destroyed by the Romans, as a pirate-haven. (Josephus, B. J. iii. 9. 3, and Strabo xvi. 2. 28.) At a later time, all describe it as an unsafe haven. Perhaps the shore changed, since the rings, to which Andromeda was tabled to have been fastened, and which probably were once used to moor vessels, were high above the sea. Perhaps, like the Channel Islands, the navigation was safe to those who knew the coast, unsafe to others. To this port Jonah "went down" from his native country, the mountain district of Zabulon. Perhaps it was not at this time in the hands of Israel. At least, the sailors were pagan. He "went down," as the man who fell among the thieves, is said to "have gone down from Jerusalem to Jericho." Luk 10:30. He "went down" from the place which God honored by His presence and protection.
And he paid the fare thereof - Jonah describes circumstantially, how he took every step to his end. He went down, found a strongly built ship going where he wished, paid his fare, embarked. He seemed now to have done all. He had severed himself from the country where his office lay. He had no further step to take. Winds and waves would do the rest. He had but to be still. He went, only to be brought back again.
"Sin brings our soul into much senselessness. For as those overtaken by heaviness of head and drunkenness, are borne on simply and at random, and, be there pit or precipice or whatever else below them, they fall into it unawares; so too, they who fall into sin, intoxicated by their desire of the object, know not what they do, see nothing before them, present or future. Tell me, Fleest thou the Lord? Wait then a little, and thou shalt learn from the event, that thou canst not escape the hands of His servant, the sea. For as soon as he embarked, it too roused its waves and raised them up on high; and as a faithful servant, finding her fellow-slave stealing some of his master's property, ceases not from giving endless trouble to those who take him in, until she recover him, so too the sea, finding and recognizing her fellow-servant, harasses the sailors unceasingly, raging, roaring, not dragging them to a tribunal but threatening to sink the vessel with all its unless they restore to her, her fellow-servant."
"The sinner "arises," because, will he, nill he, toil he must. If he shrinks from the way of God, because it is hard, he may not yet be idle. There is the way of ambition, of covetousness, of pleasure, to be trodden, which certainly are far harder. 'We wearied ourselves (Wisdom 5:7),' say the wicked, 'in the way of wickedness and destruction, yea, we have gone through deserts where there lay no way; but the way of the Lord we have not known.' Jonah would not arise, to go to Nineveh at God's command; yet he must needs arise, to flee to Tarshish from before the presence of God. What good can he have who fleeth the Good? what light, who willingly forsaketh the Light? "He goes down to Joppa." Wherever thou turnest, if thou depart from the will of God, thou goest down. Whatever glory, riches, power, honors, thou gainest, thou risest not a whit; the more thou advancest, while turned from God, the deeper and deeper thou goest down. Yet all these things are not had, without paying the price. At a price and with toil, he obtains what he desires; he receives nothing gratis, but, at great price purchases to himself storms, griefs, peril. There arises a great tempest in the sea, when various contradictory passions arise in the heart of the sinner, which take from him all tranquility and joy. There is a tempest in the sea, when God sends strong and dangerous disease, whereby the frame is in peril of being broken. There is a tempest in the sea, when, thro' rivals or competitors for the same pleasures, or the injured, or the civil magistrate, his guilt is discovered, he is laden with infamy and odium, punished, withheld from his wonted pleasures. Psa 107:23-27. "They who go down to the sea of this world, and do business in mighty waters - their soul melteth away because of trouble; they reel to and fro and stagger like a drunken man, and all their wisdom is swallowed up."
But (And) the Lord sent out - (literally 'cast along'). Jonah had done his all. Now God's part began. This He expresses by the word, "And." Jonah took "his" measures, "and" now God takes "His." He had let him have his way, as He often deals with those who rebel against Him. He lets them have their way up to a certain point. He waits, in the tranquility of His Almightiness, until they have completed their preparations; and then, when man has ended, He begins, that man may see the more that it is His doing . "He takes those who flee from Him in their flight, the wise in their counsels, sinners in their conceits and sins, and draws them back to Himself and compels them to return. Jonah thought to find rest in the sea, and lo! a tempest." Probably, God summoned back Jonah, as soon as he had completed all on his part, and sent the tempest, soon after he left the shore.
At least, such tempests often swept along that shore, and were known by their own special name, like the Euroclydon off Crete. Jonah too alone had gone down below deck to sleep, and, when the storm came, the mariners thought it possible to put back. Josephus says of that shore, "Joppa having by nature no haven, for it ends in a rough shore, mostly abrupt, but for a short space having projections, i. e., deep rocks and cliffs advancing into the sea, inclining on either side toward each other (where the traces of the chains of Andromeda yet shown accredit the antiquity of the fable,) and the north wind beating right on the shore, and dashing the high waves against the rocks which receive them, makes the station there a harborless sea. As those from Joppa were tossing here, a strong wind (called by those who sail here, the black north wind) falls upon them at daybreak, dashing straightway some of the ships against each other, some against the rocks, and some, forcing their way against the waves to the open sea, (for they fear the rocky shore ...) the breakers towering above them, sank."
The ship was like - (literally 'thought') To be broken Perhaps Jonah means by this very vivid image to exhibit the more his own dullness. He ascribes, as it were, to the ship a sense of its own danger, as she heaved and rolled and creaked and quivered under the weight of the storm which lay on her, and her masts groaned, and her yard-arms shivered. To the awakened conscience everything seems to have been alive to God's displeasure, except itself.
And cried, every man unto his God - They did what they could. "Not knowing the truth, they yet know of a Providence, and, amid religious error, know that there is an Object of reverence." In ignorance they had received one who offended God. And now God, "whom they ignorantly worshiped" Act 17:23, while they cried to the gods, who, they thought, disposed of them, heard them. They escaped with the loss of their wares, but God saved their lives and revealed Himself to them. God hears ignorant prayer, when ignorance is not willful and sin.
To lighten it of them - , literally "to lighten from against them, to lighten" what was so much "against them," what so oppressed them. "They thought that the ship was weighed down by its wonted lading, and they knew not that the whole weight was that of the fugitive prophet." "The sailors cast forth their wares," but the ship was not lightened. For the whole weight still remained, the body of the prophet, that heavy burden, not from the nature of the body, but from the burden of sin. For nothing is so onerous and heavy as sin and disobedience. Whence also Zechariah Zac 5:7 represented it under the image of lead. And David, describing its nature, said Psa 38:4, "my wickednesses are gone over my head; as a heavy burden they are too heavy for me." And Christ cried aloud to those who lived in many sins, Mat 11:28. "Come unto Me, all ye that labor and are heavy-laden, and I will refresh you."
Jonah was gone down - , probably before the beginning of the storm, not simply before the lightening of the vessel. He could hardly have fallen asleep "then." A pagan ship was a strange place for a prophet of God, not as a prophet, but as a fugitive; and so, probably, ashamed of what he had completed, he had withdrawn from sight and notice. He did not embolden himself in his sin, but shrank into himself. The conscience most commonly awakes, when the sin is done. It stands aghast as itself; but Satan, if he can, cuts off its retreat. Jonah had no retreat now, unless God had made one.
And was fast asleep - The journey to Joppa had been long and hurried; he had "fled." Sorrow and remorse completed what fatigue began. Perhaps he had given himself up to sleep, to dull his conscience. For it is said, "he lay down and was fast asleep." Grief produces sleep; from where it is said of the apostles in the night before the Lord's Passion, when Jesus "rose up from prayer and was come to His disciples, He found them sleeping for sorrow" Luk 22:45 . "Jonah slept heavily. Deep was the sleep, but it was not of pleasure but of grief; not of heartlessness, but of heavy-heartedness. For well-disposed servants soon feel their sins, as did he. For when the sin has been done, then he knows its frightfulness. For such is sin. When born, it awakens pangs in the soul which bare it, contrary to the law of our nature. For so soon as we are born, we end the travail-pangs; but sin, so soon as born, rends with pangs the thoughts which conceived it." Jonah was in a deep sleep, a sleep by which he was fast held and bound; a sleep as deep as that from which Sisera never woke. Had God allowed the ship to sink, the memory of Jonah would have been that of the fugitive prophet. As it is, his deep sleep stands as an image of the lethargy of sin . "This most deep sleep of Jonah signifies a man torpid and slumbering in error, to whom it sufficed not to flee from the face of God, but his mind, drowned in a stupor and not knowing the displeasure of God, lies asleep, steeped in security."
What meanest thou? - or rather, "what aileth thee?" (literally "what is to thee?") The shipmaster speaks of it (as it was) as a sort of disease, that he should be thus asleep in the common peril. "The shipmaster," charged, as he by office was, with the common weal of those on board, would, in the common peril, have one common prayer. It was the prophet's office to call the pagan to prayers and to calling upon God. God reproved the Scribes and Pharisees by the mouth of the children who "cried Hosanna" Mat 21:15; Jonah by the shipmaster; David by Abigail; Sa1 25:32-34; Naaman by his servants. Now too he reproves worldly priests by the devotion of laymen, sceptic intellect by the simplicity of faith.
If so be that God will think upon us - , (literally "for us") i. e., for good; as David says, Psa 40:17. "I am poor and needy, the Lord thinketh upon" (literally "for") "me." Their calling upon their own gods had failed them. Perhaps the shipmaster had seen something special about Jonah, his manner, or his prophet's garb. He does not only call Jonah's God, "thy" God, as Darius says to Daniel "thy God" Dan 6:20, but also "the God," acknowledging the God whom Jonah worshiped, to be "the God." It is not any pagan prayer which he asks Jonah to offer. It is the prayer of the creature in its need to God who can help; but knowing its own ill-desert, and the separation between itself and God, it knows not whether He will help it. So David says Psa 25:7, "Remember not the sins of my youth nor my transgressions; according to Thy mercy remember Thou me for Thy goodness' sake, O Lord."
"The shipmaster knew from experience, that it was no common storm, that the surges were an infliction borne down from God, and above human skill, and that there was no good in the master's skill. For the state of things needed another Master who ordereth the heavens, and craved the guidance from on high. So then they too left oars, sails, cables, gave their hands rest from rowing, and stretched them to heaven and called on God."
Come, and let us cast lots - Jonah too had probably prayed, and his prayers too were not heard. Probably, too, the storm had some unusual character about it, the suddenness with which it burst upon them, its violence, the quarter from where it came, its whirlwind force . "They knew the nature of the sea, and, as experienced sailors, were acquainted with the character of wind and storm, and had these waves been such as they had known before, they would never have sought by lot for the author of the threatened wreck, or, by a thing uncertain, sought to escape certain peril." God, who sent the storm to arrest Jonah and to cause him to be cast into the sea, provided that its character should set the mariners on divining, why it came. Even when working great miracles, God brings about, through man, all the forerunning events, all but the last act, in which He puts forth His might. As, in His people, he directed the lot to fall upon Achan or upon Jonathan, so here He overruled the lots of the pagan sailors to accomplish His end. " We must not, on this precedent, immediately trust in lots, or unite with this testimony that from the Acts of the Apostles, when Matthias was by lot elected to the apostolate, since the privileges of individuals cannot form a common law." "Lots," according to the ends for which they were cast, were for:
i) The lot for dividing is not wrong if not used,
1) "without any necessity, for this would be to tempt God:"
2) "if in case of necessity, not without reverence of God, as if Holy Scripture were used for an earthly end," as in determining any secular matter by opening the Bible:
3) for objects which ought to be decided otherwise, (as, an office ought to be given to the fittest:)
4) in dependence upon any other than God Pro 16:33. "The lot is cast into the lap, but the whole disposing of it is the Lord's." So then they are lawful "in secular things which cannot otherwise be conveniently distributed," or when there is no apparent reason why, in any advantage or disadvantage, one should be preferred to another." Augustine even allows that, in a time of plague or persecution, the lot might be cast to decide who should remain to administer the sacraments to the people, lest, on the one side, all should be taken away, or, on the other, the Church be deserted.
ii.) The lot for consulting, i. e., to decide what one should do, is wrong, unless in a matter of mere indifference, or under inspiration of God, or in some extreme necessity where all human means fail.
iii.) The lot for divining, i. e., to learn truth, whether of things present or future, of which we can have no human knowledge, is wrong, except by direct inspiration of God. For it is either to tempt God who has not promised so to reveal things, or, against God, to seek superhuman knowledge by ways unsanctioned by Him. Satan may readily mix himself unknown in such inquiries, as in mesmerism. Forbidden ground is his own province.
God overruled the lot in the case of Jonah, as He did the sign which the Philistines sought . "He made the heifers take the way to Bethshemesh, that the Philistines might know that the plague came to them, not by chance, but from Hilmself" . "The fugitive (Jonah) was taken by lot, not by any virtue of the lots, especially the lots of pagan, but by the will of Him who guided the uncertain lots" "The lot betrayed the culprit. Yet not even thus did they cast him over; but, even while such a tumult and storm lay on them, they held, as it were, a court in the vessel, as though in entire peace, and allowed him a hearing and defense, and sifted everything accurately, as men who were to give account of their judgment. Hear them sifting all as in a court - The roaring sea accused him; the lot convicted and witnessed against him, yet not even thus did they pronounce against him - until the accused should be the accuser of his own sin. The sailors, uneducated, untaught, imitated the good order of courts. When the sea scarcely allowed them to breathe, whence such forethought about the prophet? By the disposal of God. For God by all this instructed the prophet to be humane and mild, all but saying aloud to him; 'Imitate these uninstructed sailors. They think not lightly of one soul, nor are unsparing as to one body, thine own. But thou, for thy part, gavest up a whole city with so many myriads. They, discovering thee to be the cause of the evils which befell them, did not even thus hurry to condemn thee. Thou, having nothing whereof to accuse the Ninevites, didst sink and destroy them. Thou, when I bade thee go and by thy preaching call them to repentance, obeyedst not; these, untaught, do all, compass all, in order to recover thee, already condemned, from punishment.'"
Tell us, for whose cause - Literally "for what to whom." It may be that they thought that Jonah had been guilty toward some other. The lot had pointed him out. The mariners, still fearing to do wrong, ask him thronged questions, to know why the anger of God followed him; "what" hast thou done "to whom?" "what thine occupation?" i. e., either his ordinary occupation, whether it was displeasing to God? or this particular business in which he was engaged, and for which he had come on board. Questions so thronged have been admired in human poetry, Jerome says. For it is true to nature. They think that some one of them will draw forth the answer which they wish. It may be that they thought that his country, or people, or parents, were under the displeasure of God. But perhaps, more naturally, they wished to "know all about him," as people say. These questions must have gone home to Jonah's conscience. "What is thy business?" The office of prophet which he had left. "Whence comest thou?" From standing before God, as His minister. "What thy country? of what people art thou?" The people of God, whom he had quitted for pagan; not to win them to God, as He commanded; but, not knowing what they did, to abet him in his flight.
What is thine occupation? - They should ask themselves, who have Jonah's office to speak in the name of God, and preach repentance . "What should be thy business, who hast consecrated thyself wholly to God, whom God has loaded with daily benefits? who approachest to Him as to a Friend? "What is thy business?" To live for God, to despise the things of earth, to behold the things of heaven," to lead others heavenward.
Jonah answers simply the central point to which all these questions tended:
I am an Hebrew - This was the name by which Israel was known to foreigners. It is used in the Old Testament, only when they are spoken of by foreigners, or speak of themselves to foreigners, or when the sacred writers mention them in contrast with foreigners . So Joseph spoke of his land Gen 40:15, and the Hebrew midwives Exo 1:19, and Moses' sister Exo 2:7, and God in His commission to Moses Exo 3:18; Exo 7:16; Exo 9:1 as to Pharaoh, and Moses in fulfilling it Exo 5:3. They had the name, as having passed the River Euphrates, "emigrants." The title might serve to remind themselves, that they were "strangers" and "pilgrims," Heb 11:13. whose fathers had left their home at God's command and for God , "passers by, through this world to death, and through death to immortality."
And I fear the Lord - , i. e., I am a worshiper of Him, most commonly, one who habitually stands in awe of Him, and so one who stands in awe of sin too. For none really fear God, none fear Him as sons, who do not fear Him in act. To be afraid of God is not to fear Him. To be afraid of God keeps men away from God; to fear God draws them to Him. Here, however, Jonah probably meant to tell them, that the Object of his fear and worship was the One Self-existing God, He who alone is, who made all things, in whose hands are all things. He had told them before, that he had fled "from being before Yahweh." They had not thought anything of this, for they thought of Yahweh, only as the God of the Jews. Now he adds, that He, Whose service he had thus forsaken, was "the God of heaven, Who made the sea and dry land," that sea, whose raging terrified them and threatened their lives. The title, "the God of heaven," asserts the doctrine of the creation of the heavens by God, and His supremacy.
Hence, Abraham uses it to his servant Gen 24:7, and Jonah to the pagan mariners, and Daniel to Nebuchadnezzar Dan 2:37, Dan 2:44; and Cyrus in acknowledging God in his proclamation Ch2 36:23; Ezr 1:2. After his example, it is used in the decrees of Darius Ezr 6:9-10 and Artaxerxes Ezr 7:12, Ezr 7:21, Ezr 7:23, and the returned exiles use it in giving account of their building the temple to the Governor Ezr 5:11-12. Perhaps, from the habit of contact with the pagan, it is used once by Daniel Dan 2:18 and by Nehemiah Neh 1:4-5; Neh 2:4, Neh 2:20. Melchizedek, not perhaps being acquainted with the special name, Yahweh, blessed Abraham in the name of "God, the Possessor" or "Creator of heaven and earth" Gen 14:19, i. e., of all that is. Jonah, by using it, at once taught the sailors that there is One Lord of all, and why this evil had fallen on them, because they had himself with them, the renegade servant of God. "When Jonah said this, he indeed feared God and repented of his sin. If he lost filial fear by fleeing and disobeying, he recovered it by repentance."
Then were the men exceedingly afraid - Before, they had feared the tempest and the loss of their lives. Now they feared God. They feared, not the creature but the Creator. They knew that what they had feared was the doing of His Almightiness. They felt how awesome a thing it was to be in His Hands. Such fear is the beginning of conversion, when people turn from dwelling on the distresses which surround them, to God who sent them.
Why hast thou done this? - They are words of amazement and wonder. Why hast thou not obeyed so great a God, and how thoughtest thou to escape the hand of the Creator ? "What is the mystery of thy flight? Why did one, who feared God and had revelations from God, flee, sooner than go to fulfill them? Why did the worshiper of the One true God depart from his God?" "A servant flee from his Lord, a son from his father, man from his God!" The inconsistency of believers is the marvel of the young Christian, the repulsion of those without, the hardening of the unbeliever. If people really believed in eternity, how could they be thus immersed in things of time? If they believed in hell, how could they so hurry there? If they believed that God died for them, how could they so requite Him? Faith without love, knowledge without obedience, conscious dependence and rebellion, to be favored by God yet to despise His favor, are the strangest marvels of this mysterious world.
All nature seems to cry out to and against the unfaithful Christian, "why hast thou done this?" And what a why it is! A scoffer has recently said so truthfully : "Avowed scepticism cannot do a tenth part of the injury to practical faith, that the constant spectacle of the huge mass of worldly unreal belief does." It is nothing strange, that the world or unsanctified intellect should reject the Gospel. It is a thing of course, unless it be converted. But, to know, to believe, and to DISOBEY! To disobey God, in the name of God. To propose to halve the living Gospel, as the woman who had killed her child Kg1 3:26, and to think that the poor quivering remnants would be the living Gospel anymore! As though the will of God might, like those lower forms of His animal creation, be divided endlessly, and, keep what fragments we will, it would still be a living whole, a vessel of His Spirit! Such unrealities and inconsistencies would be a sore trial of faith, had not Jesus, who (cf. Joh 2:25), "knew what is in man," forewarned us that it should be so. The scandals against the Gospel, so contrary to all human opinion, are only all the more a testimony to the divine knowledge of the Redeemer.
What shall we do unto thee? - They knew him to be a prophet; they ask him the mind of his God. The lots had marked out Jonah as the cause of the storm; Jonah had himself admitted it, and that the storm was for "his" cause, and came from "his" God . "Great was he who fled, greater He who required him. They dare not give him up; they cannot conceal him. They blame the fault; they confess their fear; they ask "him" the remedy, who was the author of the sin. If it was faulty to receive thee, what can we do, that God should not be angered? It is thine to direct; ours, to obey."
The sea wrought and was tempestuous - , literally "was going and whirling." It was not only increasingly tempestuous, but, like a thing alive and obeying its Master's will, it was holding on its course, its wild waves tossing themselves, and marching on like battalions, marshalled, arrayed for the end for which they were sent, pursuing and demanding the runaway slave of God . "It was going, as it was bidden; it was going to avenge its Lord; it was going, pursuing the fugitive prophet. It was swelling every moment, and, as though the sailors were too tardy, was rising in yet greater surges, shewing that the vengeance of the Creator admitted not of delay."
Take me up, and cast me into the sea - Neither might Jonah have said this, nor might the sailors have obeyed it, without the command of God. Jonah might will alone to perish, who had alone offended; but, without the command of God, the Giver of life, neither Jonah nor the sailors might dispose of the life of Jonah. But God willed that Jonah should be cast into the sea - where he had gone for refuge - that (Wisdom 11:16) wherewithal he had "sinned, by the same also he might be punished" as a man; and, as a prophet, that he might, in his three days' burial, prefigure Him who, after His Resurrection, should convert, not Nineveh, but the world, the cry of whose wickedness went up to God.
For I know that for my sake - o "In that he says, "I know," he marks that he had a revelation; in that he says, "this great storm," he marks the need which lay on those who cast him into the sea."
The men rowed hard - , literally "dug." The word, like our "plowed the main," describes the great efforts which they made. Amid the violence of the storm, they had furled their sails. These were worse than useless. The wind was off shore, since by rowing alpine they hoped to get back to it. They put their oars well and firmly in the sea, and turned up the water, as men turn up earth by digging. But in vain! God willed it not. The sea went on its way, as before. In the description of the deluge, it is repeated Gen 7:17-18, "the waters increased and bare up the ark, and it was lifted up above the earth; the waters increased greatly upon the earth; and the ark went upon the face of the waters." The waters raged and swelled, drowned the whole world, yet only bore up the ark, as a steed bears its rider: man was still, the waters obeyed. In this tempest, on the contrary, man strove, but, instead of the peace of the ark, the burden is, the violence of the tempest; "the sea wrought and was tempestuous against them" . "The prophet had pronounced sentence against himself, but they would not lay hands upon him, striving hard to get back to land, and escape the risk of bloodshed, willing to lose life rather than cause its loss. O what a change was there. The people who had served God, said, Crucify Him, Crucify Him! These are bidden to put to death; the sea rageth; the tempest commandeth; and they are careless its to their own safety, while anxious about another's."
Wherefore (And) they cried unto the Lord - "They cried" no more "each man to his god," but to the one God, whom Jonah had made known to them; and to Him they cried with an earnest submissive, cry, repeating the words of beseeching, as men, do in great earnestness; "we beseech Thee, O Lord, let us not, we beseech Thee, perish for the life of this man" (i. e., as a penalty for taking it, as it is said, Sa2 14:7. "we will slay him for the life of his brother," and, Deu 19:21. "life for life.") They seem to have known what is said, Gen 9:5-6. "your blood of your lives will I require; at the hand of every beast will I require it and at the hand of man; at the hand of every man's brother will I require the life of man. Whoso sheddeth man's blood, by man shall his blood be shed, for in the image of God made He man" , "Do not these words of the sailors seem to us to be the confession of Pilate, who washed his hands, and said, 'I am clean from the blood of this Man?' The Gentiles would not that Christ should perish; they protest that His Blood is innocent."
And lay not upon us innocent blood - innocent as to them, although, as to this thing, guilty before God, and yet, as to God also, more innocent, they would think, than they. For, strange as this was, one disobedience, their whole life, they now knew, was disobedience to God; His life was but one act in a life of obedience. If God so punishes one sin of the holy Pe1 4:18, "where shall the ungodly and sinner appear?" Terrible to the awakened conscience are God's chastenings on some (as it seems) single offence of those whom He loves.
For Thou, Lord, (Who knowest the hearts of all men,) hast done, as it pleased Thee - Wonderful, concise, confession of faith in these new converts! Psalmists said it, Psa 135:6; Psa 115:3. "Whatsoever God willeth, that doeth He in heaven and in earth, in the sea and in all deep places." But these had but just known God, and they resolve the whole mystery of man's agency and God's Providence into the three simple words , as (Thou) "willedst" (Thou) "didst." "That we took him aboard, that the storm ariseth, that the winds rage, that the billows lift themselves, that the fugitive is betrayed by the lot, that he points out what is to be done, it is of Thy will, O Lord" . "The tempest itself speaketh, that 'Thou, Lord, hast done as Thou willedst.' Thy will is fulfilled by our hands." "Observe the counsel of God, that, of his own will, not by violence or by necessity, should he be cast into the sea. For the casting of Jonah into the sea signified the entrance of Christ into the bitterness of the Passion, which He took upon Himself of His own will, not of necessity. Isa 53:7. "He was offered up, and He willingly submitted Himself." And as those who sailed with Jonah were delivered, so the faithful in the Passion of Christ. Joh 18:8-9. "If ye seek Me, let these go their way, that the saying might be fulfilled which" Jesus spake, 'Of them which Thou gavest Me, I have lost none. '"
They took up Jonah - o "He does not say, 'laid hold on him', nor 'came upon him' but 'lifted' him; as it were, bearing him with respect and honor, they cast him into the sea, not resisting, but yielding himself to their will."
The sea ceased (literally "stood") from his raging - Ordinarily, the waves still swell, when the wind has ceased. The sea, when it had received Jonah, was hushed at once, to show that God alone raised and quelled it. It "stood" still, like a servant, when it had accomplished its mission. God, who at all times saith to it Job 38:11, "Hitherto shalt thou come and no further, and here shall thy proud waves be stayed," now unseen, as afterward in the flesh Mat 8:26, "rebuked the winds and the sea, and there was a great calm" . "If we consider the errors of the world before the Passion of Christ, and the conflicting blasts of diverse doctrines, and the vessel, and the whole race of man, i. e., the creature of the Lord, imperiled, and, after His Passion, the tranquility of faith and the peace of the world and the security of all things and the conversion to God, we shall see how, after Jonah was cast in, the sea stood from its raging" . "Jonah, in the sea, a fugitive, shipwrecked, dead, sayeth the tempest-tossed vessel; he sayeth the pagan, aforetime tossed to and fro by the error of the world into divers opinions. And Hosea, Amos, Isaiah, Joel, who prophesied at the same time, could not amend the people in Judaea; whence it appeared that the breakers could not be calmed, save by the death of (Him typified by) the fugitive."
And the men feared the Lord with a great fear - because, from the tranquility of the sea and the ceasing of the tempest, they saw that the prophet's words were true. This great miracle completed the conversion of the mariners. God had removed all human cause of fear; and yet, in the same words as before, he says, "they feared a great fear;" but he adds, "the Lord." It was the great fear, with which even the disciples of Jesus feared, when they saw the miracles which He did, which made even Peter say, Luk 5:8. "Depart from me, for I am a sinful man, O Lord." Events full of wonder had thronged upon them; things beyond nature, and contrary to nature; tidings which betokened His presence, Who had all things in His hands. They had seen "wind and storm fulfilling His word" Psa 148:8, and, forerunners of the fishermen of Galilee, knowing full well from their own experience that this was above nature, they felt a great awe of God. So He commanded His people, "Thou shalt fear the Lord thy God Deu 6:13, for thy good always" Deu 6:24.
And offered a sacrifice - Doubtless, as it was a large decked vessel and bound on a long voyage, they had live creatures on board, which they could offer in sacrifice. But this was not enough for their thankfulness; "they vowed vows." They promised that they would do thereafter what they could not do then ; "that they would never depart from Him whom they had begun to worship." This was true love, not to be content with aught which they could do, but to stretch forward in thought to an abiding and enlarged obedience, as God should enable them. And so they were doubtless enrolled among the people of God, firstfruits from among the pagan, won to God Who overrules all things, through the disobedience and repentance of His prophet. Perhaps, they were the first preachers among the pagan, and their account of their own wonderful deliverance prepared the way for Jonah's mission to Nineveh.
Now the Lord had (literally "And the Lord") prepared - Jonah (as appears from his thanksgiving) was not swallowed at once, but sank to the bottom of the sea, God preserving him in life there by miracle, as he did in the fish's belly. Then, when the seaweed was twined around his head, and he seemed to be already buried until the sea should give up her dead, "God prepared the fish to swallow Jonah" . "God could as easily have kept Jonah alive in the sea as in the fish's belly, but, in order to prefigure the burial of the Lord, He willed him to be within the fish whose belly was as a grave." Jonah, does not say what fish it was; and our Lord too used a name, signifying only one of the very largest fish. Yet it was no greater miracle to create a fish which should swallow Jonah, than to preserve him alive when swallowed . "The infant is buried, as it were, in the womb of its mother; it cannot breathe, and yet, thus too, it liveth and is preserved, wonderfully nurtured by the will of God." He who preserves the embryo in its living grave can maintain the life of man as easily without the outward air as with it.
The same Divine Will preserves in being the whole creation, or creates it. The same will of God keeps us in life by breathing this outward air, which preserved Jonah without it. How long will men think of God, as if He were man, of the Creator as if He were a creature, as though creation were but one intricate piece of machinery, which is to go on, ringing its regular changes until it shall be worn out, and God were shut up, as a sort of mainspring within it, who might be allowed to be a primal Force, to set it in motion, but must not be allowed to vary what He has once made? "We must admit of the agency of God," say these men when they would not in name be atheists, "once in the beginning of things, but must allow of His interference as sparingly as may be." Most wise arrangement of the creature, if it were indeed the god of its God! Most considerate provision for the non-interference of its Maker, if it could but secure that He would not interfere with it for ever! Acute physical philosophy, which, by its omnipotent word, would undo the acts of God! Heartless, senseless, sightless world, which exists in God, is upheld by God, whose every breath is an effluence of God's love, and which yet sees Him not, thanks Him not, thinks it a greater thing to hold its own frail existence from some imagined law, than to be the object of the tender personal care of the Infinite God who is Love! Poor hoodwinked souls, which would extinguish for themselves the Light of the world, in order that it may not eclipse the rushlight of their own theory!
And Jonah was in the belly of the fish - The time that Jonah was in the fish's belly was a hidden prophecy. Jonah does not explain nor point it. He tells the fact, as Scripture is accustomed to do so. Then he singles out one, the turning point in it. Doubtless in those three days and nights of darkness, Jonah (like him who after his conversion became Paul), meditated much, repented much, sorrowed much, for the love of God, that he had ever offended God, purposed future obedience, adored God with wondering awe for His judgment and mercy. It was a narrow home, in which Jonah, by miracle, was not consumed; by miracle, breathed; by miracle, retained his senses in that fetid place. Jonah doubtless, repented, marveled, adored, loved God. But, of all, God has singled out this one point, how, out of such a place, Jonah thanked God. As He delivered Paul and Silas from the prison, when they prayed with a loud voice to Him, so when Jonah, by inspiration of His Spirit, thanked Him, He delivered him.
To thank God, only in order to obtain fresh gifts from Him, would be but a refined, hypocritical form of selfishness. Such a formal act would not be thanks at all. We thank God, because we love Him, because He is so infinitely good, and so good to us, unworthy. Thanklessness shuts the door to His personal mercies to us, because it makes them the occasion of fresh sins of our's. Thankfulness sets God's essential goodness free (so to speak) to be good to us. He can do what He delights in doing, be good to us, without our making His Goodness a source of harm to us. Thanking Him through His grace, we become fit vessels for larger graces . "Blessed he who, at every gift of grace, returns to Him in whom is all fullness of graces; to whom when we show ourselves not ungrateful for gifts received, we make room in ourselves for grace, and become meet for receiving yet more." But Jonah's was that special character of thankfulness, which thanks God in the midst of calamities from which there was no human exit; and God set His seal on this sort of thankfulness, by annexing this deliverance, which has consecrated Jonah as an image of our Lord, to his wonderful act of thanksgiving. |
|Met proto-oncogene (hepatocyte growth factor receptor)|
Crystallographic structure of MET. PDB rendering based on 1r0p.
|External IDs||ChEMBL: GeneCards:|
|RNA expression pattern|
c-Met (MET or MNNG HOS Transforming gene) is a proto-oncogene that encodes a protein known as hepatocyte growth factor receptor (HGFR). The hepatocyte growth factor receptor protein possesses tyrosine-kinase activity. The primary single chain precursor protein is post-translationally cleaved to produce the alpha and beta subunits, which are disulfide linked to form the mature receptor.
MET is a membrane receptor that is essential for embryonic development and wound healing. Hepatocyte growth factor (HGF) is the only known ligand of the MET receptor. MET is normally expressed by cells of epithelial origin, while expression of HGF is restricted to cells of mesenchymal origin. Upon HGF stimulation, MET induces several biological responses that collectively give rise to a program known as invasive growth.
Abnormal MET activation in cancer correlates with poor prognosis, where aberrantly active MET triggers tumor growth, formation of new blood vessels (angiogenesis) that supply the tumor with nutrients, and cancer spread to other organs (metastasis). MET is deregulated in many types of human malignancies, including cancers of kidney, liver, stomach, breast, and brain. Normally, only stem cells and progenitor cells express MET, which allows these cells to grow invasively in order to generate new tissues in an embryo or regenerate damaged tissues in an adult. However, cancer stem cells are thought to hijack the ability of normal stem cells to express MET, and thus become the cause of cancer persistence and spread to other sites in the body.
MET proto-oncogene (GeneID: 4233) has a total length of 125,982 bp, and it is located in the 7q31 locus of chromosome 7. MET is transcribed into a 6,641 bp mature mRNA, which is then translated into a 1,390 amino-acid MET protein.
MET is a receptor tyrosine kinase (RTK) that is produced as a single-chain precursor. The precursor is proteolytically cleaved at a furin site to yield a highly glycosylated extracellular α-subunit and a transmembrane β-subunit, which are linked together by a disulfide bridge.
- Region of homology to semaphorins (Sema domain), which includes the full α-chain and the N-terminal part of the β-chain
- Cysteine-rich MET-related sequence (MRS domain)
- Glycine-proline-rich repeats (G-P repeats)
- Four immunoglobulin-like structures (Ig domains), a typical protein-protein interaction region.
A Juxtamembrane segment that contains:
- a serine residue (Ser 985), which inhibits the receptor kinase activity upon phosphorylation
- a tyrosine (Tyr 1003), which is responsible for MET polyubiquitination, endocytosis, and degradation upon interaction with the ubiquitin ligase CBL
- Tyrosine kinase domain, which mediates MET biological activity. Following MET activation, transphosphorylation occurs on Tyr 1234 and Tyr 1235
- C-terminal region contains two crucial tyrosines (Tyr 1349 and Tyr 1356), which are inserted into the multisubstrate docking site, capable of recruiting downstream adapter proteins with Src homology-2 (SH2) domains. The two tyrosines of the docking site have been reported to be necessary and sufficient for the signal transduction both in vitro.
MET signaling pathway
MET activation by its ligand HGF induces MET kinase catalytic activity, which triggers transphosphorylation of the tyrosines Tyr 1234 and Tyr 1235. These two tyrosines engage various signal transducers, thus initiating a whole spectrum of biological activities driven by MET, collectively known as the invasive growth program. The transducers interact with the intracellular multisubstrate docking site of MET either directly, such as GRB2, SHC, SRC, and the p85 regulatory subunit of phosphatidylinositol-3 kinase (PI3K), or indirectly through the scaffolding protein Gab1
Tyr 1349 and Tyr 1356 of the multisubstrate docking site are both involved in the interaction with GAB1, SRC, and SHC, while only Tyr 1356 is involved in the recruitment of GRB2, phospholipase C γ (PLC-γ), p85, and SHP2.
GAB1 is a key coordinator of the cellular responses to MET and binds the MET intracellular region with high avidity, but low affinity. Upon interaction with MET, GAB1 becomes phosphorylated on several tyrosine residues which, in turn, recruit a number of signalling effectors, including PI3K, SHP2, and PLC-γ. GAB1 phosphorylation by MET results in a sustained signal that mediates most of the downstream signaling pathways.
Activation of signal transduction
MET engagement activates multiple signal transduction pathways:
- The RAS pathway mediates HGF-induced scattering and proliferation signals, which lead to branching morphogenesis. Of note, HGF, differently from most mitogens, induces sustained RAS activation, and thus prolonged MAPK activity.
- The PI3K pathway is activated in two ways: PI3K can be either downstream of RAS, or it can be recruited directly through the multifunctional docking site. Activation of the PI3K pathway is currently associated with cell motility through remodeling of adhesion to the extracellular matrix as well as localized recruitment of transducers involved in cytoskeletal reorganization, such as RAC1 and PAK. PI3K activation also triggers a survival signal due to activation of the AKT pathway.
- The STAT pathway, together with the sustained MAPK activation, is necessary for the HGF-induced branching morphogenesis. MET activates the STAT3 transcription factor directly, through an SH2 domain.
- The beta-catenin pathway, a key component of the Wnt signaling pathway, translocates into the nucleus following MET activation and participates in transcriptional regulation of numerous genes.
Role in development
During embryonic development, transformation of the flat, two-layer germinal disc into a three-dimensional body depends on transition of some cells from an epithelial phenotype to spindle-shaped cells with motile behaviour, a mesenchymal phenotype. This process is referred to as epithelial-mesenchymal transition (EMT). Later in embryonic development, MET is crucial for gastrulation, angiogenesis, myoblast migration, bone remodeling, and nerve sprouting among others. MET is essential for embryogenesis, because MET -/- mice die in utero due to severe defects in placental development. Furthermore, MET is required for such critical processes as liver regeneration and wound healing during adulthood.
Tissue distribution
MET is normally expressed by epithelial cells. However, MET is also found on endothelial cells, neurons, hepatocytes, hematopoietic cells, and melanocytes. HGF expression is restricted to cells of mesenchymal origin.
Transcriptional control
MET transcription is activated by HGF and several growth factors. MET promoter has four putative binding sites for Ets, a family of transcription factors that control several invasive growth genes. ETS1 activates MET transcription in vitro. MET transcription is activated by hypoxia-inducible factor 1 (HIF1), which is activated by low concentration of intracellular oxygen. HIF1 can bind to one of the several hypoxia response elements (HREs) in the MET promoter. Hypoxia also activates transcription factor AP-1, which is involved in MET transcription.
Role in cancer
MET pathway plays an important role in the development of cancer through:
- angiogenesis (sprouting of new blood vessels from pre-existing ones to supply a tumor with nutrients);
- scatter (cells dissociation due to metalloprotease production), which often leads to metastasis.
Coordinated down-regulation of both MET and its downstream effector extracellular signal-regulated kinase 2 (ERK2) by miR-199a* may be effective in inhibiting not only cell proliferation but also motility and invasive capabilities of tumor cells.
Interaction with tumour suppressor genes
PTEN (phosphatase and tensin homolog) is a tumor suppressor gene encoding a protein PTEN, which possesses lipid and protein phosphatase-dependent as well as phosphatase-independent activities. PTEN protein phosphatase is able to interfere with MET signaling by dephosphorylating either PIP3 generated by PI3K, or the p52 isoform of SHC. SHC dephosphorylation inhibits recruitment of the GRB2 adapter to activated MET.
Cancer therapies targeting HGF/MET
Since tumor invasion and metastasis are the main cause of death in cancer patients, interfering with MET signaling appears to be a promising therapeutic approach. A comprehensive list of HGF and MET targeted experimental therapeutics for oncology now in human clinical trials can be found here.
MET kinase inhibitors
Kinase inhibitors are low molecular weight molecules that prevent ATP binding to MET, thus inhibiting receptor transphosphorylation and recruitment of the downstream effectors. The limitations of kinase inhibitors include the facts that they only inhibit kinase-dependent MET activation, and that none of them is fully specific for MET.
- K252a (Fermentek Biotechnology) is a staurosporine analogue isolated from Nocardiopsis sp. soil fungi, and it is a potent inhibitor of all receptor tyrosine kinases (RTKs). At nanomolar concentrations, K252a inhibits both the wild type and the mutant (M1268T) MET function.
- SU11274 (SUGEN) specifically inhibits MET kinase activity and its subsequent signaling. SU11274 is also an effective inhibitor of the M1268T and H1112Y MET mutants, but not the L1213V and Y1248H mutants. SU11274 has been demonstrated to inhibit HGF-induced motility and invasion of epithelial and carcinoma cells.
- PHA-665752 (Pfizer) specifically inhibits MET kinase activity, and it has been demonstrated to represses both HGF-dependent and constitutive MET phosphorylation. Furthermore, some tumors harboring MET amplifications are highly sensitive to treatment with PHA-665752.
- ARQ197 (ArQule) is a promising selective inhibitor of MET, which entered a phase 2 clinical trial in 2008.
- Foretinib (XL880, Exelixis) targets multiple receptor tyrosine kinases (RTKs) with growth-promoting and angiogenic properties. The primary targets of foretinib are MET, VEGFR2, and KDR. Foretinib has completed a phase 2 clinical trials with indications for papillary renal cell carcinoma, gastric cancer, and head and neck cancer.
- SGX523 (SGX Pharmaceuticals) specifically inhibits MET at low nanomolar concentrations.
- MP470 (SuperGen) is a novel inhibitor of c-KIT, MET, PDGFR, Flt3, and AXL. Phase I clinical trial of MP470 had been announced in 2007.
HGF inhibitors
Since HGF is the only known ligand of MET, formation of a HGF:MET complex blocks MET biological activity. For this purpose, truncated HGF, anti-HGF neutralizing antibodies, and an uncleavable form of HGF have been utilized so far. The major limitation of HGF inhibitors is that they block only HGF-dependent MET activation.
- NK4 competes with HGF as it binds MET without inducing receptor activation, thus behaving as a full antagonist. NK4 is a molecule bearing the N-terminal hairpin and the four kringle domains of HGF. Moreover, NK4 is structurally similar to angiostatins, which is why it possesses anti-angiogenic activity.
- Neutralizing anti-HGF antibodies were initially tested in combination, and it was shown that at least three antibodies, acting on different HGF epitopes, are necessary to prevent MET tyrosine kinase activation. More recently, it has been demonstrated that fully human monoclonal antibodies can individually bind and neutralize human HGF, leading to regression of tumors in mouse models. Two anti-HGF antibodies are currently available: the humanized AV299 (AVEO), and the fully human AMG102 (Amgen).
- Uncleavable HGF is an engineered form of pro-HGF carrying a single amino-acid substitution, which prevents the maturation of the molecule. Uncleavable HGF is capable of blocking MET-induced biological responses by binding MET with high affinity and displacing mature HGF. Moreover, uncleavable HGF competes with the wild-type endogenous pro-HGF for the catalytic domain of proteases that cleave HGF precursors. Local and systemic expression of uncleavable HGF inhibits tumor growth and, more importantly, prevents metastasis.
Decoy MET
Decoy MET refers to a soluble truncated MET receptor. Decoys are able to inhibit MET activation mediated by both HGF-dependent and independent mechanisms, as decoys prevent both the ligand binding and the MET receptor homodimerization. CGEN241 (Compugen) is a decoy MET that is highly efficient in inhibiting tumor growth and preventing metastasis in animal models.
Immunotherapy targeting MET
Drugs used for immunotherapy can act either passively by enhancing the immunologic response to MET-expressing tumor cells, or actively by stimulating immune cells and altering differentiation/growth of tumor cells.
Passive immunotherapy
Administering monoclonal antibodies (mAbs) is a form of passive immunotherapy. MAbs facilitate destruction of tumor cells by complement-dependent cytotoxicity (CDC) and cell-mediated cytotoxicity (ADCC). In CDC, mAbs bind to specific antigen, leading to activation of the complement cascade, which in turn leads to formation of pores in tumor cells. In ADCC, the Fab domain of a mAb binds to a tumor antigen, and Fc domain binds to Fc receptors present on effector cells (phagocytes and NK cells), thus forming a bridge between an effector and a target cells. This induces the effector cell activation, leading to phagocytosis of the tumor cell by neutrophils and macrophages. Furthermore, NK cells release cytotoxic molecules, which lyse tumor cells.
- DN30 is monoclonal anti-MET antibody that recognizes the extracellular portion of MET. DN30 induces both shedding of the MET ectodomain as well as cleavage of the intracellular domain, which is successively degraded by proteasome machinery. As a consequence, on one side MET is inactivated, and on the other side the shed portion of extracellular MET hampers activation of other MET receptors, acting as a decoy. DN30 inhibits tumour growth and prevents metastasis in animal models.
- OA-5D5 is one-armed monoclonal anti-MET antibody that was demonstrated to inhibit orthotopic pancreatic and glioblastoma tumor growth and to improve survival in tumor xenograft models. OA-5D5 is produced as a recombinant protein in Escherichia coli. It is composed of murine variable domains for the heavy and light chains with human IgG1 constant domains. The antibody blocks HGF binding to MET in a competitive fashion.
Active immunotherapy
Active immunotherapy to MET-expressing tumors can be achieved by administering cytokines, such as interferons (IFNs) and interleukins (IL-2), which triggers non-specific stimulation of numerous immune cells. IFNs have been tested as therapies for many types of cancers and have demonstrated therapeutic benefits. IL-2 has been approved by the U.S. Food and Drug Administration (FDA) for the treatment of renal cell carcinoma and metastatic melanoma, which often have deregulated MET activity.
See also
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Further reading
- Peruzzi B, Bottaro DP (2006). "Targeting the c-Met signaling pathway in cancer". Clin. Cancer Res. 12 (12): 3657–60. doi:10.1158/1078-0432.CCR-06-0818. PMID 16778093.
- Birchmeier C, Birchmeier W, Gherardi E, Vande Woude GF (December 2003). "Met, metastasis, motility and more". Nat. Rev. Mol. Cell Biol. 4 (12): 915–25. doi:10.1038/nrm1261. PMID 14685170.
- Zhang YW, Vande Woude GF (February 2003). "HGF/SF-met signaling in the control of branching morphogenesis and invasion". J. Cell. Biochem. 88 (2): 408–17. doi:10.1002/jcb.10358. PMID 12520544.
- Paumelle R, Tulasne D, Kherrouche Z, Plaza S, Leroy C, Reveneau S, Vandenbunder B, Fafeur V, Tulashe D, Reveneau S (April 2002). "Hepatocyte growth factor/scatter factor activates the ETS1 transcription factor by a RAS-RAF-MEK-ERK signaling pathway". Oncogene 21 (15): 2309–19. doi:10.1038/sj.onc.1205297. PMID 11948414.
- Comoglio PM (1993). "Structure, biosynthesis and biochemical properties of the HGF receptor in normal and malignant cells". EXS 65: 131–65. PMID 8380735.
- Maulik G, Shrikhande A, Kijima T, et al. (2002). "Role of the hepatocyte growth factor receptor, c-Met, in oncogenesis and potential for therapeutic inhibition". Cytokine Growth Factor Rev. 13 (1): 41–59. doi:10.1016/S1359-6101(01)00029-6. PMID 11750879.
- Ma PC, Maulik G, Christensen J, Salgia R (2004). "c-Met: structure, functions and potential for therapeutic inhibition". Cancer Metastasis Rev. 22 (4): 309–25. doi:10.1023/A:1023768811842. PMID 12884908.
- Knudsen BS, Edlund M (2004). "Prostate cancer and the met hepatocyte growth factor receptor". Adv. Cancer Res. Advances in Cancer Research 91: 31–67. doi:10.1016/S0065-230X(04)91002-0. ISBN 978-0-12-006691-9. PMID 15327888.
- Dharmawardana PG, Giubellino A, Bottaro DP (2005). "Hereditary papillary renal carcinoma type I". Curr. Mol. Med. 4 (8): 855–68. doi:10.2174/1566524043359674. PMID 15579033.
- Kemp LE, Mulloy B, Gherardi E (2006). "Signalling by HGF/SF and Met: the role of heparan sulphate co-receptors". Biochem. Soc. Trans. 34 (Pt 3): 414–7. doi:10.1042/BST0340414. PMID 16709175.
- Proto-Oncogene Proteins c-met at the US National Library of Medicine Medical Subject Headings (MeSH)
- UniProtKB/Swiss-Prot entry P08581: MET_HUMAN, ExPASy (Expert Protein Analysis System) proteomics server of the Swiss Institute of Bioinformatics (SIB)
- A table with references to significant roles of MET in cancer |
Copyright (c) Arvin S. Quist
INTRODUCTION TO CLASSIFICATION
THE NEED FOR CLASSIFICATION
A government is responsible for the survival of the nation and its people. To ensure that survival, a government must sometimes stringently control certain information that (1) gives the nation a significant advantage over adversaries or (2) prevents adversaries from having an advantage that could significantly damage the nation. Governments protect that special information by classifying it; that is, by giving it a special designation, such as "Secret," and then restricting access to it (e.g., by need-to-know requirements and physical security measures).
This right of a government to keep certain information concerning national security (secrets) from most of the nation's citizens is nearly universally accepted. Since antiquity, governments have protected information that gave them an advantage over adversaries. In wartime, when a nation's survival is at stake, the reasons for secrecy are most apparent, the secrecy restrictions imposed by the government are most widespread,[*] and acceptance of those restrictions by the citizens is broadest.[†] In peacetime, there are fewer reasons for secrecy in government, generally the government classifies less information, and citizens are less willing to accept security restrictions on information.
MAJOR AREAS OF CLASSIFIED INFORMATION
The information that is classified by most democracies, whether in peacetime or wartime, is usually limited to information that concerns the nation's defense or its foreign relations--military and diplomatic information. Most of that information falls within five major areas: (1) military operations, (2) weapons technology, (3) diplomatic activities, (4) intelligence activities, and (5) cryptology. The latter two areas might be considered to be special parts of the first three areas. That is, intelligence and cryptology are "service" functions for the primary areas--military operations, weapons technology, and diplomatic activities. From a historical perspective, the classification of weapons technology became widespread only in the 20th century. Classification of information about military operations and diplomatic activities has been practiced for millennia.
Examples of military-operations information that is frequently classified include information concerning the strength and deployment of forces, troop movements, ship sailings, the location and timing of planned attacks, tactics and strategy, and supply logistics. Obviously, if an enemy learned the major details of an impending attack, that attack would be less successful than if it came as a surprise to the enemy.* Information possessed by a government about an adversary's military activities or capabilities must be protected to preserve the ability to predict those activities or to neutralize those capabilities. If the adversary knew that the government had this information, the adversary would change those plans or capabilities. Military-operations information is usually classified for only a limited time. After an operation is over, most of the important information is known to the enemy.
Weapons technology is classified to preserve the advantage of surprise in the first use of a new weapon,† to prevent an adversary from developing effective countermeasures against a new weapon,‡ or to prevent an adversary from using that technology against its originator (by developing a similar weapon). A major factor in that latter reason for classifying weapons technology is "lead time." Classifying advanced weapons-technology information prevents an adversary from using that information to shorten the time required to produce similar weapons systems for its own use. Consequently, assuming continued advancements in a weapons technology by the initial developer of that technology, the adversary's weapons systems will not be as effective as those of the nation that initially developed that technology, and the adversary will be at a disadvantage.
With respect to lead time, when weapons systems can be significantly improved, then information on "obsolete" weapons is much less sensitive than information on newer weapons. Thus, information on muzzle-loading rifle technology was not as sensitive as that on breech-loading rifle technology, which was not as sensitive as information on lever-action rifle technology, . . . semiautomatic rifle . . . automatic rifle . . . machine gun. However, with respect to nuclear weapons, a "rogue" nation or terrorist group can probably achieve its objectives just as easily with "crude" kiloton nuclear weapons that might require a ship or truck to transport as with sophisticated megaton nuclear weapons that might fit into a (large) suitcase. Thus, "obsolete" nuclear-weapons technology should be continue to be protected, especially with respect to technologies concerning production of highly enriched uranium or other nuclear-weapon materials.
Weapons technology includes scientific and technical information related to that technology. World War I marked the start of the "modern" period when science and technology affected the development of weapons systems to a greater degree than any time previously. That interrelationship became even more pronounced in World War II, with notable scientific and technological successes: the atomic bomb, radar, and the proximity fuse. World War II, particularly with respect to the atomic bomb, marked the first time that the progress of military technology was significantly influenced by scientists, as contrasted to advances by engineers or by scientists working as engineers.
With respect to classification, the more that applied scientific or technical information is uniquely applicable to weapons, the more likely that this information will be classified. Generally, basic research is not classified unless it represents a major breakthrough leading to a completely new weapons system. An example of that circumstance was the rigid classification during World War II, and for several years thereafter, of much basic scientific research related to atomic energy (nuclear weapons).
The need for secrecy in diplomatic negotiations and relations has long been recognized. A nation's ability to obtain favorable terms in negotiations with other countries would be diminished if its negotiating strategy and goals were known in advance to the other countries.* The effectiveness of military-assistance agreements between nations would be impaired if an adversary knew of them and could plan to neutralize them. In New York Times v. United States, the "Pentagon Papers" case, U.S. Supreme Court Justice Stewart recognized the importance of secrecy in foreign policy and national defense matters:
It is elementary that the successful conduct of international diplomacy and the maintenance of an effective national defense requires both confidentiality and secrecy. Other nations can hardly deal with this Nation in an atmosphere of mutual trust unless they know that their confidences will be kept . . .. In the area of basic national defense the frequent need for absolute secrecy is, of course, self evident.
During the term of the first
president, it was established that some need for secrecy in diplomatic matters would remain even after negotiations were completed. President Washington, in 1796, refused a request by the House of Representatives for documents prepared for treaty negotiations with U.S. and gave the following as one reason for refusal: England
The nature of foreign negotiations requires caution, and their success must often depend on secrecy; and even when brought to a conclusion a full disclosure of all the measures, demands, or eventual concessions which may have been proposed or contemplated would be extremely impolitic; for this might have a pernicious influence on future negotiations, or produce immediate inconvenience, perhaps danger and mischief, in relation to other powers.
It has been said that President Nixon initially was not going to attempt to stop the New York Times and other newspapers from publishing the "Pentagon Papers." However, the executive branch was then in secret diplomatic negotiations with
, and Henry Kissinger "is said to have persuaded the president that the Chinese wouldn't continue their secret parleys if they saw that China couldn't keep its secrets." Washington
Intelligence information includes information gathering and covert operations. Collecting military and diplomatic information about other nations involves the use of photoreconnaissance airplanes and satellites, communication intercepts, the review of documents obtained openly, and other overt methods. However, information gathering also includes the use of undercover agents, confidential sources, and other covert methods. For those covert activities, secrecy is usually imposed on the identity of agents or sources, on information about intelligence methods and capabilities, and on much of the information received from the covert sources. Few clandestine agents could be recruited (or, in some instances, would live long) if their identity were not a closely guarded secret. Information provided by a clandestine agent must frequently be classified because, if a government knew that some of its information was compromised, it might be able to determine the identity of the person (agent) who provided the information to its adversary. Successful intelligence-gathering methods must be protected so that the adversary does not know the degree of their success and is not stimulated to develop countermeasures to stop the flow of information. Intelligence information from friendly nations is generally classified by the recipient country. Allies would be less willing to share intelligence information if they knew that it would not be protected against disclosure.
Cryptology encompasses methods to code and transmit secret messages and methods to intercept and decode messages. Writing messages in code, or cryptography,* has been practiced for thousands of years. One of the earliest preserved texts of a coded message is an inscription carved on an Egyptian tomb in about 1900 B.C. The earliest known pottery glaze formula was written in code on a Mesopotamian cuneiform tablet in about 1500 B.C. The Spartans established a system of military cryptography by the 5th century B.C. Persia later used cryptography for political purposes. Cryptography began its steady development in western civilization starting about the 13th century, primarily in
. By the early 16th century, Italy 's ruling Council of Ten had an elaborate organization for enciphering and deciphering messages. Venice
Restrictions on cryptologic information are necessary to protect
communications. Diplomatic negotiations could not successfully be conducted at locations other than the seat of government if safe communications could not be established. Cryptologic information must also be protected to prevent an adversary from learning of a nation's capabilities to intercept and decode messages. If an adversary learns that its communications are not secure, it will use another method, which will require additional time and effort to defeat.[‡] The Allies' World War II success in breaking the German codes contributed to shortening that war. That success was kept secret until 1974, about 34 years after the German code had been broken and about 29 years after World War II had ended. The U.S. Army's success in breaking a World War II U.S.S.R. code (the Venona project, which began in 1943 and continued until 1980) was not made public until about 1995. That was about 50 years after the first such message had been deciphered (and about 45 years after the U.S.S.R. had learned through espionage of the Army's success). U.S.
BASIS FOR CLASSIFICATION IN THE UNITED STATES
The need for governmental secrecy was directly recognized in the U.S. Constitution. Article I, Sect. 5, of the Constitution explicitly authorizes secrecy in government by stating that "Each House shall keep a Journal of its Proceedings, and from time to time publish the same, excepting such Parts as in their Judgment require Secrecy." Also included in the Constitution, in Article I, Sect. 9, is a statement that "a regular Statement and Account of the Receipts and Expenditures of all public Money shall be published from time to time." A U.S. Court of Appeals has determined that the phrase "from time to time" was intended to authorize expenditures for certain military or foreign relations matters that were intended to be kept secret for a time.
The Constitution does not explicitly provide for secrecy by the Executive Branch of the U.S. Government. However, the authority of that Executive Branch to keep certain information secret from most
citizens is implicit in its executive responsibilities, which include the national defense and foreign relations. This presidential authority has been upheld by the Supreme Court in a number of cases. Judicial decisions have also relied on a common-law privilege for a government to withhold information concerning national defense and foreign relations. Congress, by two statutes, the Freedom of Information Act and the Internal Security Act of 1950, has implicitly recognized the president's authority to classify information (see Chapter 3). U.S.
At this time in the
, information is classified either by presidential authority, currently Executive Order 12958, or by statute, the Atomic Energy Act of 1954, as amended (Atomic Energy Act). Classification under Executive Orders and under the Atomic Energy Act is extensively discussed in Chapters 3 and 4, respectively. United States
CLASSIFICATION AND SECURITY
Classification has been variously described as the "cornerstone" of national security, the "mother" of security, and the "kingpin" of an information security system.,,, Classification identifies the information that must be protected against unauthorized disclosure. Security determines how to protect information after it is classified. Security includes both personnel security and physical security.
The initial classification determination, establishing what should not be disclosed to adversaries and the level of protection required, is probably the most important single factor in the security of all classified projects and programs., None of the expensive personnel-clearance and information-control provisions (physical security aspects) of an information security system comes into effect until information has been classified; classification is the pivot on which the whole subsequent security system turns (excluding security for other reasons, such as to prevent theft of materials). 19 Therefore, it is important to classify only information that truly warrants protection in the interest of national security.
Since the mid 1970s, several classification experts have remarked on the increasing emphasis by some government agencies on physical-security matters, which has been accompanied by a decreased emphasis on the classification function. One of the founders (and the first chairman) of the National Classification Management Society (NCMS), who was also an Atomic Energy Commission Contractor Classification Officer, has expressed concern about the tendency to emphasize the word "security" at the expense of the word "classification" with respect to security classification of information.17 In the mid 1980s another charter member of the NCMS pointed out that, although the status of classification still remained high in the Department of Energy (DOE), the situation had changed within the Department of Defense, where Classification Management had been organizationally placed under Security. Even the NCMS, founded as a classification organization, appears to be changing to become increasingly oriented towards security matters rather than classification matters. It is noteworthy that the marked emphasis by the U.S. Government in recent years on physical-security measures has not been accompanied by any significant increased emphasis on classification matters.
The previous paragraph was written in 1989, and the trend described in that paragraph has continued. The classification function at DOE headquarters is now a part of the security organization as is the classification function at many DOE operations offices and DOE-contractor organizations. That function generally used to be part of a technical or other non-security organization. The NCMS has also continued to become more security-oriented.
With respect to classification as a profession (or lack of recognition thereof), it is interesting to note some comments and a recommendation in the Report of the Commission on Protecting and Reducing Government Secrecy. In this 1997 report, that Commission noted the "all-important initial decision of whether to classify at all," and that "this first step of the classification management process . . . tends to be the weakest link in the process of identifying, marking, and then protecting the information." The Commission further stated that "the importance of the initial decision to classify cannot be overstated." However, the Commission then stated that "classification and declassification policy and oversight . . . should be viewed primarily as information management issues which require personnel with subject matter and records management expertise." Although recommending that "The Federal Government . . . [should] create, support, and promote an information systems security career field within the Government," the Commission made no similar recommendation for security classification of information as a profession or career. Res ipsa loquitur.
[*] "When a nation is at war many things that might be said in time of peace are such a hindrance to its effort that their utterance will not be endured so long as men fight and that no Court could regard them as protected by any constitutional right" [Schenck v. United States, 249 U.S. 47, 52 (1919) (J. Holmes)].
[†] Since the September 11, 2001, terrorist attacks against the World Trade Center towers and the Pentagon, the United States considers itself to be in a war against terrorism. One consequence has been a significant shift in opinion, not only of the general public but also of some strong supporters of freedom-of-information matters, towards favoring more control of information that might aid terrorists. This increased control, especially pertaining to weapons of mass destruction, includes (1) establishing broader criteria for identifying information that is classified or "sensitive"; (2) permitting reclassification of declassified information, and (3) restricting further governmental distribution of documents already released to the public.
*However, during the Greek and Roman eras in the Mediterranean, when the infantry was paramount and both sides were approximately equally equipped with respect to weapons, many battles were fought without attempts to maintain secrecy of troop movements or with respect to surprise attacks (B. and F. M. Brodie, From Crossbow to H-Bomb, Indiana University Press, Bloomington, Ind., 1973, p. 17).
†"Secret" weapons have proven decisive in warfare. One example of the decisive impact of a new weapon was at the battle of Crecy in 1346. At this battle, the English used their "secret" weapon, the longbow, to defeat the French decisively. Although the French had a two-to-one superiority in numbers (about 40,000 to 20,000), the French lost about 11,500 men, while the English lost only about 100 men (W. S. Churchill, A History of the English-Speaking Peoples, Vol. 1, Dodd, Mead and Co., New York, 1961, pp. 332-351; B. and F. M. Brodie, From Crossbow to H-Bomb, Indiana University Press, Bloomington, Ind., 1973, pp. 37-40).
‡In World War II, the Germans developed an acoustic torpedo designed to home in on a ship's propellers. However, the Allies obtained advance information about this torpedo so that when it was first used by the Germans, countermeasures were already in place (B. and F. M. Brodie, From Crossbows to H-Bombs, Indiana University Press, Bloomington, Ind., 1973, p. 222).
*In 1921, the United States, Britain, France, Italy, and Japan held a conference to limit their naval armaments. The United States had broken Japan's diplomatic code and thereby knew the lowest naval armaments that Japan would accept. Therefore, U.S. negotiators had merely to wait out Japan's negotiators to reach terms favorable to the United States (J. Bamford, The Puzzle Palace, Houghton, Mifflin Co., Boston, 1982, pp. 9-10).
*The breaking of codes is termed cryptanalysis.
[‡] Even "friendly" nations get upset if they know that one of their codes has been broken. As noted earlier in this chapter, the United States deciphered Japan's diplomatic code in 1921. Herbert O. Yardley, who was principally responsible for breaking this code, wrote a book, The American Black Chamber, published in 1931, which included information on this matter. Yardley's book did not contribute to developing friendly United States-Japanese relations. A consequence of this revelation was enactment of a U.S. statute that made it a crime for anyone who, by virtue of his employment by the United States, obtained access to a diplomatic code or a message in such code and published or furnished to another such code or message, "or any matter which was obtained while in the process of transmission between any foreign government and its diplomatic mission in the United States" (48 Stat. 122, June 10, 1933, codified at 18 U.S.C. Sect. 952.)
B. and F. M. Brodie, From Crossbow to H-Bomb, Indiana University Press, Bloomington, Ind., 1973, p. 172. Hereafter this book is cited as "Brodie."
Brodie, p. 233.
New York Times v. United States, 403 U.S. 713, 728 (1971).
J. D. Richardson, A Compilation of Messages and Papers of the Presidents. 1789-1897, U.S. Government Printing Office, Washington, D.C., Vol. I, at 194-195 (1896).
Richard Gid Powers, "Introduction," in Secrecy--The American Experience, by Daniel Patrick Moynihan, Yale University Press, New Haven, Conn., 1998, p. 32.
D. Kahn, The Codebreakers, MacMillan, Inc., New York, 1967, p. 71. Hereafter cited as "Kahn."
Kahn, p. 75.
Kahn, p. 82.
Kahn, p. 86.
Kahn, p. 106.
Kahn, p. 109.
See, for example, F. W. Winterbotham, The Ultra Secret, Harper & Row, New York, 1974.
Halperin v. CIA, 629 F.2d 144, 154-162 (D.C. Cir., 1980).
U.S. Constitution, Article II, sect. 2.
See, for example, Totten v. United States, 92 U.S. 105 (1875); United States v. Reynolds, 345 U.S. 1 (1952); Weinberger v. Catholic Action of Hawaii, 454 U.S. 139 (1981).
F. E. Rourke, Secrecy and Publicity: Dilemmas of Democracy, Johns Hopkins Press, Baltimore, 1961, pp. 63-64.
D. B. Woodbridge, "Footnotes," J. Natl. Class. Mgmt. Soc. 12 (2), 120-124 (1977), p.122.
R. J. Boberg, "Panel--Classification Management Today," J. Natl. Class. Mgmt. Soc. 5 (2), 56-60 (1969), p. 57.
E. J. Suto, "History of Classification," J. Natl. Class. Mgmt. Soc. 12 (1), 9-17 (1976), p.13.
James J. Bagley, "NCMS - Now and the Future," J. Natl. Class. Mgmt. Soc. 25, 20-29 (1989), p. 28.
T. S. Church, "Panel--Science and Technology, and Classification Management," J. Natl. Class. Mgmt. Soc. 2, 39-45 (1966), p. 40.
W. N. Thompson, "Security Classification Management Coordination Between Industry and DOD," J. Natl. Class. Mgmt. Soc. 4 (2), 121-128 (1969), p. 121.
W. N. Thompson, "User Agency Security Classification Management and Program Security," J. Natl. Class. Mgmt. Soc. 8, 52-53 (1972), p. 52.
Department of Defense Handbook for Writing Security Classification Guidance, DoD 5200.1-H, U.S. Department of Defense, Mar. 1986, p. 1-1.
F. J. Daigle, "Woodbridge Award Acceptance Remarks," J. Natl. Class. Mgmt. Soc. 21, 110-112 (1985), p. 111.
D. C. Richardson, "Management or Enforcement," J. Natl. Class. Mgmt. Soc. 23, 13-20 (1987).
Report of the Commission on Protecting and Reducing Government Secrecy, S. Doc. 105-2, Daniel Patrick Moynihan, Chairman; Larry Combest, Vice Chairman, Commission on Protecting and Reducing Government Secrecy, U.S. Government Printing Office, Washington, D.C., 1997. Hereafter cited as the "Moynihan Report."
Moynihan Report, p. 19.
Moynihan Report, p. 35.
Moynihan Report, p. 44.
Moynihan Report, p. 111. |
On April 30, 1789, George Washington was sworn in as the first President of the United States.
George Washington wrote the following on the eve of his inauguration:
It is said that every man has his portion of ambition. I may have mine, I suppose, as well as the rest, but if I know my own heart, my ambition would not lead me into public life; my only ambition is to do my duty in this world as well as I am capable of performing it, and to merit the good opinion of all good men.
We are so lucky, so very lucky, to have had this man in our “canon”. There’s as always, so much to say. One of the thing that strikes me about him is that he never wanted to seem like he was jostling for power or position. George Washington had many wonderful qualities and abilities – but it was this distaste for public life that I believe made him truly great. He went out of his way to let everyone know how unworthy he felt, how he hoped their trust in him was warranted, that he was eager to finally go home and live the life of a private man… But on this day in history, April 30, there was to be no private man anymore. His people had chosen him, and while Mount Vernon continued to call to him, he knew he must accept.
David McCullough describes, in his book on John Adams, inauguration day:
On the day of his inauguration, Thursday, April 30 1789, Washington rode to Federal Hall in a canary-yellow carriage pulled by six white horses and followed by a long column of New York militia in full dress. The air was sharp, the sun shone brightly, and with all work stopped in the city, the crowds along his route were the largest ever seen. It was as if all New York had turned out and more besides. “Many persons in the crowd,” reported the Gazette of the United States “were heard to say they should now die contented � nothing being wanted to complete their happiness � but the sight of the savior of his country.”
In the Senate Chamber were gathered the members of both houses of Congress, the Vice President, and sundry officials and diplomatic agents, all of whom rose when Washington made his entrance, looking solemn and stately. His hair powdered, he wore a dress sword, white silk stockings, shoes with silver buckles, and a suit of the same brown Hartford broadcloth that Adams, too, was wearing for the occasion. They might have been dressed as twins, except that Washington’s metal buttons had eagles on them.
It was Adams who formally welcomed the General and escorted him to the dais. For an awkward moment Adams appeared to be in some difficulty, as though he had forgotten what he was supposed to say. then, addressing Washington, he declared that the Senate and House of Representatives were ready to attend him for the oath of office as required by the Constitution. Washington said he was ready. Adams bowed and led the way to the outer balcony, in full view of the throng in the streets. People were cheering and waving from below, and from windows and rooftops as far as the eye could see. Washington bowed once, then a second time.
Fourteen years earlier, it had been Adams who called on the Continental Congress to make the tall Virginian commander-in-chief of the army. Now he stood at Washington’s side as Washington, his right hand on the Bible, repeated the oath of office as read by Chancellor Robert R. Livingston of New York, who had also been a member of the Continental Congress.
In a low voice Washington solemnly swore to execute the office of the President of the United States and, to the best of his ability, to “preserve, protect, and defend the Constitution of the United States.” Then, as not specified in the Constitution, he added, “So help me God”, and kissed the Bible, thereby establishing his own first presidential tradition.
“It is done,” Livingston said, and, turning to the crowd, cried out, “Long live George Washington, President of the United States.”
The following is George Washington’s first inaugural address. What I sense in these words is what I sense in so many of the original documents of that time, written by the main players: they were embarking on a grand and hopeful experiment. They were entering uncharted waters. And they all seem determined (each in their different ways, with their different views) to make the most of the opportunity, to seize the day. No decision was unimportant, everything had meaning … and what I also sense in this inaugural address is that Washington knew that he wasn’t only talking to the people present, but he was also talking to us. The future generations. They all knew that they were being watched, carefully, by those who would come after.
The only thing required of a President on his inauguration day, in those early early days, was that he take the oath of Office. Washington, in composing an address, to the people who put their faith in him, set the precedent. Every president since then has followed his example.
George Washington’s first inaugural address:
Fellow-Citizens of the Senate and of the House of Representatives:
Among the vicissitudes incident to life no event could have filled me with greater anxieties than that of which the notification was transmitted by your order, and received on the 14th day of the present month. On the one hand, I was summoned by my Country, whose voice I can never hear but with veneration and love, from a retreat which I had chosen with the fondest predilection, and, in my flattering hopes, with an immutable decision, as the asylum of my declining years–a retreat which was rendered every day more necessary as well as more dear to me by the addition of habit to inclination, and of frequent interruptions in my health to the gradual waste committed on it by time. On the other hand, the magnitude and difficulty of the trust to which the voice of my country called me, being sufficient to awaken in the wisest and most experienced of her citizens a distrustful scrutiny into his qualifications, could not but overwhelm with despondence one who (inheriting inferior endowments from nature and unpracticed in the duties of civil administration) ought to be peculiarly conscious of his own deficiencies. In this conflict of emotions all I dare aver is that it has been my faithful study to collect my duty from a just appreciation of every circumstance by which it might be affected. All I dare hope is that if, in executing this task, I have been too much swayed by a grateful remembrance of former instances, or by an affectionate sensibility to this transcendent proof of the confidence of my fellow-citizens, and have thence too little consulted my incapacity as well as disinclination for the weighty and untried cares before me, my error will be palliated by the motives which mislead me, and its consequences be judged by my country with some share of the partiality in which they originated.
Such being the impressions under which I have, in obedience to the public summons, repaired to the present station, it would be peculiarly improper to omit in this first official act my fervent supplications to that Almighty Being who rules over the universe, who presides in the councils of nations, and whose providential aids can supply every human defect, that His benediction may consecrate to the liberties and happiness of the people of the United States a Government instituted by themselves for these essential purposes, and may enable every instrument employed in its administration to execute with success the functions allotted to his charge. In tendering this homage to the Great Author of every public and private good, I assure myself that it expresses your sentiments not less than my own, nor those of my fellow- citizens at large less than either. No people can be bound to acknowledge and adore the Invisible Hand which conducts the affairs of men more than those of the United States. Every step by which they have advanced to the character of an independent nation seems to have been distinguished by some token of providential agency; and in the important revolution just accomplished in the system of their united government the tranquil deliberations and voluntary consent of so many distinct communities from which the event has resulted can not be compared with the means by which most governments have been established without some return of pious gratitude, along with an humble anticipation of the future blessings which the past seem to presage. These reflections, arising out of the present crisis, have forced themselves too strongly on my mind to be suppressed. You will join with me, I trust, in thinking that there are none under the influence of which the proceedings of a new and free government can more auspiciously commence.
By the article establishing the executive department it is made the duty of the President “to recommend to your consideration such measures as he shall judge necessary and expedient.” The circumstances under which I now meet you will acquit me from entering into that subject further than to refer to the great constitutional charter under which you are assembled, and which, in defining your powers, designates the objects to which your attention is to be given. It will be more consistent with those circumstances, and far more congenial with the feelings which actuate me, to substitute, in place of a recommendation of particular measures, the tribute that is due to the talents, the rectitude, and the patriotism which adorn the characters selected to devise and adopt them. In these honorable qualifications I behold the surest pledges that as on one side no local prejudices or attachments, no separate views nor party animosities, will misdirect the comprehensive and equal eye which ought to watch over this great assemblage of communities and interests, so, on another, that the foundation of our national policy will be laid in the pure and immutable principles of private morality, and the preeminence of free government be exemplified by all the attributes which can win the affections of its citizens and command the respect of the world. I dwell on this prospect with every satisfaction which an ardent love for my country can inspire, since there is no truth more thoroughly established than that there exists in the economy and course of nature an indissoluble union between virtue and happiness; between duty and advantage; between the genuine maxims of an honest and magnanimous policy and the solid rewards of public prosperity and felicity; since we ought to be no less persuaded that the propitious smiles of Heaven can never be expected on a nation that disregards the eternal rules of order and right which Heaven itself has ordained; and since the preservation of the sacred fire of liberty and the destiny of the republican model of government are justly considered, perhaps, as deeply, as finally, staked on the experiment entrusted to the hands of the American people.
Besides the ordinary objects submitted to your care, it will remain with your judgment to decide how far an exercise of the occasional power delegated by the fifth article of the Constitution is rendered expedient at the present juncture by the nature of objections which have been urged against the system, or by the degree of inquietude which has given birth to them. Instead of undertaking particular recommendations on this subject, in which I could be guided by no lights derived from official opportunities, I shall again give way to my entire confidence in your discernment and pursuit of the public good; for I assure myself that whilst you carefully avoid every alteration which might endanger the benefits of an united and effective government, or which ought to await the future lessons of experience, a reverence for the characteristic rights of freemen and a regard for the public harmony will sufficiently influence your deliberations on the question how far the former can be impregnably fortified or the latter be safely and advantageously promoted.
To the foregoing observations I have one to add, which will be most properly addressed to the House of Representatives. It concerns myself, and will therefore be as brief as possible. When I was first honored with a call into the service of my country, then on the eve of an arduous struggle for its liberties, the light in which I contemplated my duty required that I should renounce every pecuniary compensation. From this resolution I have in no instance departed; and being still under the impressions which produced it, I must decline as inapplicable to myself any share in the personal emoluments which may be indispensably included in a permanent provision for the executive department, and must accordingly pray that the pecuniary estimates for the station in which I am placed may during my continuance in it be limited to such actual expenditures as the public good may be thought to require.
Having thus imparted to you my sentiments as they have been awakened by the occasion which brings us together, I shall take my present leave; but not without resorting once more to the benign Parent of the Human Race in humble supplication that, since He has been pleased to favor the American people with opportunities for deliberating in perfect tranquillity, and dispositions for deciding with unparalleled unanimity on a form of government for the security of their union and the advancement of their happiness, so His divine blessing may be equally conspicuous in the enlarged views, the temperate consultations, and the wise measures on which the success of this Government must depend.
William Maclay, a senator from Pennsylvania, kept a daily journal – highly detailed, and rather cynical, about the Senate sessions of the first Congress. He describes the first inauguration in vivid detail:
30th April, Thursday.–This is a great, important day. Goddess of etiquette, assist me while I describe it. The Senate stood adjourned to half after eleven o’clock. About ten dressed in my best clothes; went for Mr. Morris’ lodgings, but met his son, who told me that his father would not be in town until Saturday. Turned into the Hall. The crowd already great. The Senate met. The Vice-President rose in the most solemn manner. This son of Adam seemed impressed with deeper gravity, yet what shall I think of him? He often, in the midst of his most important airs–I believe when tie is at loss for expressions (and this he often is, wrapped up, I suppose, in the contemplation of his own importance)– suffers an unmeaning kind of vacant laugh to escape him. This was the case to-day, and really to me bore the air of ridiculing the farce he was acting. “Gentlemen, I wish for the direction of the Senate. The President will, I suppose, addressthe Congress. How shall I behave? How shall we receive it? Shall it be standing or sitting?”
Here followed a considerable deal of talk from him which I could make nothing of. Mr. Lee began with the House of Commons (as is usual with him), then the House of Lords, then the King, and then back again. The result of his information was, that the Lords sat and the Commons stood on the delivery of the King’s speech. Mr. Izard got up and told how often he had been in the Houses of Parliament. He said a great deal of what he had seen there. [He] made, however, this sagacious discovery, that the Commons stood because they had no. seats to sit on, being arrived at the bar of the House of Lords. It was discovered after some time that the King sat, too, and had his robes and crown on.
Mr. Adams got up again and said he had been very often indeed at the Parliament on those occasions, but there always was such a crowd, and ladies along, that for his part he could not say how it was. Mr. Carrol got up to declare that he thought it of no consequence how it was in Great Britain; they were no rule to us, etc. But all at once the Secretary, who had been out, whispered to the Chair that the Clerk from the Representatives was at the door with a communication. Gentlemen of the Senate, how shall he be received? A silly kind of resolution of the committee on that business had been laid on the table some days ago. The amount of it was that each House should communicate to the other what and how they chose; it concluded, however, something in this way: That everything should be done with all the propriety that was proper. The question was, Shall this be adopted, that we may know how to receive the Clerk? It was objected [that] this will throw no light on the subject; it will leave you where you are. Mr. Lee brought the House of Commons before us again. He reprobated the rule; declared that the Clerk should not come within the bar of file House; that the proper mode was for the Sergeant-at-Arms, with the mace on his shoulder, to meet the Clerk at the door and receive his communication; we are not, however, provided for this ceremonious way of doing business, having neither mace nor sergeant nor Masters in Chancery, who carry down bills from the English Lords.
Mr. Izard got up and labored unintelligibly to show the great distinction between a communication and a delivery of a thing, but he was not minded. Mr. Elsworth showed plainly enough that if the Clerk was not permitted to deliver the communication, the Speaker might as well send it inclosed. Repeated accounts came [that] the Speaker and Representatives were at the door. Confusion ensued; the members left their seats. Mr. Read rose and called the attention of the Senate to the neglect that had been shown Mr. Thompson, late Secretary. Mr. Lee rose to answer him, but I could not hear one word he said. The Speaker was introduced, followed by the Representatives. Here we sat an hour and ten minutes before the President arrived–this delay was owing to Lee, Izard, and Dalton, who had stayed with us while the Speaker came in, instead of going to attend the President. The President advanced between the Senate and Representatives, bowing to each. He was placed in the chair by the Vice-President; the Senate with their president on the right, the Speaker and the Representatives on his left. The Vice-President rose and addressed a short sentence to him. The import of it was that he should now take the oath of office as President. He seemed to have forgot half what he was to say, for he made a dead pause and stood for some time, to appearance, in a vacant mood. He finished with a formal bow, and the President was conducted out of the middle window into the gallery, and the oath was administered by the Chancellor. Notice that the business done was communicated to the crowd by proclamation, etc., who gave three cheers, and repeated it on the President’s bowing to them.
As the company returned into the Senate chamber, the President took the chair and the Senators and Representatives their seats. He rose, and all arose also and addressed them. This great man was agitated and embarrassed more than ever he was by the leveled cannon or pointed musket. He trembled, and several times could scarce make out to read, though it must be supposed he had often read it before. He put part of the fingers of his left hand into the side of what I think the tailors call the fall of the breeches, changing the paper into his left hand. After some time he then did the same with some of the fingers of his right hand. When he came to the words all the world, he made a flourish with his right hand, which left rather an ungainly impression. I sincerely, for my part, wished all set ceremony in the hands of the dancing-masters, and that this first of men had read off his address in the plainest manner, without ever taking his eyes from the paper, for I felt hurt that he was not first in everything. He was dressed in deep brown, with metal buttons, with an eagle on them, white stockings, a bag, and sword.
From the hall there was a grand procession to Saint Paul’s Church, where prayers were said by the Bishop. The procession was well conducted and without accident, as far as I have heard. The militia were all under arms, lined the street near the church, made a good figure, and behaved well.
The Senate returned to their chamber after service, formed, and took up the address. Our Vice-President called it his most gracious speech. I can not approve of this. A committee was appointed on it–Johnson, Carrol, Patterson. Adjourned. In the evening there were grand fireworks. The Spanish Ambassador’s house was adorned with transparent paintings; the French Minister’s house was illuminated, and had some transparent pieces; the Hall was grandly illuminated, and after all this the people went to bed.
I have such a deep fondness for John Adams, with all his airs and self-importance and vanity. I just love the guy, what can I say. He’s so feckin’ human.
The description of Washington’s awkwardness makes me want to cry:
He rose, and all arose also and addressed them. This great man was agitated and embarrassed more than ever he was by the leveled cannon or pointed musket. He trembled, and several times could scarce make out to read, though it must be supposed he had often read it before. He put part of the fingers of his left hand into the side of what I think the tailors call the fall of the breeches, changing the paper into his left hand. After some time he then did the same with some of the fingers of his right hand. When he came to the words all the world, he made a flourish with his right hand, which left rather an ungainly impression. I sincerely, for my part, wished all set ceremony in the hands of the dancing-masters, and that this first of men had read off his address in the plainest manner, without ever taking his eyes from the paper, for I felt hurt that he was not first in everything. He was dressed in deep brown, with metal buttons, with an eagle on them, white stockings, a bag, and sword.
God. Good God. But what really moves me is that after the address, they all walked in procession, led by George Washington, to St. Paul’s Church, for a service.
St. Paul’s Church. (Read that article … it’s a well-known story, of course, but it always bears repeating.) St. Paul’s has always had meaning for us here in New York, because of its long history, but now … it has more meaning than ever. I can’t even think about St. Paul’s without feeling tears come to my eyes. So to think … that that special church, that church that became symbolic (not just to us here, but to people all over the country) of hope, or survival, of healing … would be the place where George Washington prayed for guidance after being sworn in as the first President… I mean, honestly. I don’t even know what else to say about it.
April 30, 1789 … the day this new nation embarked on its unknown and exciting course, with George Washington at the helm.
Here is an image of the first page of this inaugural address, in Washington’s own hand. |
- Definition of syncopation in the Online Dictionary. Meaning of syncopation. Pronunciation of syncopation. Translations of syncopation. syncopation synonyms, syncopation antonyms. Information about syncopation in the free online English. — “syncopation - definition of syncopation by the Free Online”,
- Also, if the musician suddenly does not play anything on beat 1, that would also be syncopation. Playing a note ever-so-slightly before or after a beat is another form of syncopation because this produces an unexpected accent. — “Syncopation - Definition”,
- Syncopation. In music, syncopation includes a variety of rhythms which are in some way unexpected in that they deviate from the strict succession of regularly spaced strong and weak but also powerful beats in a meter (pulse) In music, syncopation includes a variety of rhythms which are in some. — “Syncopation”,
- Syncopation definition, a shifting of the normal accent, usually by stressing the normally unaccented beats. See more. — “Syncopation | Define Syncopation at ”,
- Syncopations can happen anywhere: in the melody, the bass line, the rhythm section, the chordal accompaniment. Ragtime, for example, would hardly be ragtime without the jaunty syncopations in the melody set against the steady unsyncopated bass. — “Syncopation”,
- Syncopation Software is a leading provider of decision support software tools including decision ***ysis software, risk ***ysis Founded as a spin-off from a first-tier consulting organization, Syncopation develops products that go beyond the basics to give you exceptional insight into real-world. — “DPL Decision Tree & Business Risk ***ysis Software”,
- More simply, syncopation is a general term for a disturbance or interruption of the Syncopation is used in many musical styles, and is fundamental in black-influenced styles. — “Syncopation - Simple English Wikipedia, the free encyclopedia”,
- (Click to enlarge) syncopation from Mozart's Symphony no. 25 syncopation n. Music . A shift of accent in a passage or composition that occurs when a. — “syncopation: Definition from ”,
- Welcome to ! Practicing Smart • Musical Theory • Musical Syncopation • Improvising on the Piano • Laurel Webster. Copyright2010 Loraine & Chris Plante. All Rights Reserved. website design by C. — “Musical Syncopation”,
- Syncopation.TV is the means by which we exhibit, ***yze and discuss the fascinating world of music across genre, cultural, age and national boundaries. This is the place where you can further explore music and the men and women dedicated to this art and craft. — “Syncopation.TV - Home”, syncopation.tv
- Syncopation ( pronounced SINK-o-PAY-shun) is hearing the beat when you don't expect to. When notes begin before or after a strong beat, you have syncopation. Listen to this rhythm in Example 1. The notes fall on the beats. This is not syncopation. Now listen to Example 2. — “Syncopation”, empire.k12.ca.us
- I've been testing out Syncopation using the new iTunes release from Apple and have found Syncopation 2.2 to be fully compatible with iTunes 10. Syncopation has been thoroughly tested on Snow Leopard, and I'm happy to say that no compatibility issues have been found. — “Sonzea - Home”,
- In music the word syncopation has a very specific meaning. Actually, the main syncopation in cha-cha happens on beat two (usually the first part of a rock step), I'll explain why shortly. — “Syncopation in dance and music”,
- on | off. enter. — “Syncopation Official Web Site”,
- For other uses of the same name, see Syncopation (disambiguation). In music, syncopation includes a variety of rhythms which are in some way unexpected in that they deviate from the strict succession of regularly spaced strong and weak but also powerful beats in a meter (pulse). — “Syncopation - Wikipedia, the free encyclopedia”,
- Syncopation is a musical process that involves adding an unexpected element to the basic beat of a musical composition. At times, the syncopation adds more beats, while at other times it delays or changes the sense of a particular beat in the line of rhythm. — “What Is Syncopation?”,
- Dynamic Syncopation music profile on Yahoo! Music. Find lyrics, free streaming MP3s, music videos and photos of Dynamic Syncopation on Yahoo! Music. — “Dynamic Syncopation on Yahoo! Music”,
- Syncopation Makes It Move. Introduction. Syncopation is the wonderful effect which is created when notes which are supposed to fall in a certain place surprise the ear and fall somewhere else. When the accent falls off the beat, or if there is silence on the beat, then syncopation occurs. — “Music Concepts Online: Syncopation Lecture”, neiu.edu
- Syncopation. Syncopation is a disturbance or interruption of the regular flow of rhythm. Syncopation shifts this emphasis, or, to put it another way, it places the accent on the wrong syllable. — “Syncopation: Excerpts from The NPR Classical Music Companion”, kennedy-
- syncopation in a folk context, check out the examples below. In this example, there is a syncopated melody over a simple quarter-note bass part, which. — “Syncopation”,
- Syncopation is used in many musical styles, including classical music, but it is fundamental in such styles as reggae, ragtime, rap, jump blues, jazz and often in dubstep. In the form of a back beat, syncopation is used in virtually all contemporary popular music. — “Syncopation - New World Encyclopedia”,
- The Elite Syncopation quintet has been touring the United States from coast to coast for over a decade, Elite Syncopation tours nationally, giving dozens of concerts a year, both locally in Connecticut and throughout the rest of the country. Click here for booking information. — “Elite Syncopation: Who We Are”,
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“Syncopation November 14, 2008. In Knitting, My designs | 1 comment. Bet you didn't think I'd actually get this done today, did you. Pattern: Syncopation, of my own devising. Model: Rowanspun DK. Yup, this yarn is discontinued. I have it on good authority”
— Syncopation | Fiber Dreams,
“Guitar Noize is a guitar blog focussing on guitar news, crazy guitar designs, Syncopation experiment”
— Syncopation experiment,
“Forum Intro. Greeley Central Drumline Forum. http://syncopation.aimoo. and Exercises. Schedules. Sectionals >> Go to Greeley Central Drumline Forum. Similar Forums”
— Greeley Central Drumline Forum - Aimoo,
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“Before starting to explain what musical syncopation is, it is response, or trackback from your own site. Posted in Music by Mickys-Blog No Comments Yet. Leave a Comment”
— What is Musical Syncopation And How Does It Work?,
“Syncopation. Saturday, 27 October 2007. If you are Pagan and have a blog, then we would To join Syncopation's Pagan Blog List: You must display one of the buttons which are on”
“Songwriting and Music Forum > Community Blog > Songstuff Stuff > syncopation. off. Subscribe to Songstuff Our resident drumming expert Tom Hoffman introduces Syncopation and discusses why it is so important for drummers”
— Songwriting and Music Forum -> Songstuff Stuff,
“Submitted by syncopation on Thu, 03/29/2007 - 9:43pm. Related topics: Submitted by syncopation on Thu, 03/29/2007 - 5:31pm. Related topics: love. poem. poetry”
— syncopation's blog | The Icarus Project, |
I love the fall and how the leaves change from deep greens to reds and orange and gold. This natural riot of color takes place wherever there are trees with leaves and there’s almost no place better to watch the leaves change than in the Northeast. This part of the four-seasoned ritual of life attracts tourists from far and wide and tugs at me to make a special trip to our home in the mountains there. And this reminds me every year about the natural changes that are a constant in our lives.
Ever wonder why and how the leaves change colors?
• As summer ends and autumn comes, the days get shorter and shorter. This is how the trees "know" to begin getting ready for winter. The trees will begin to rest and live off the food they stored during the summer. The green chlorophyll disappears from the leaves. As the bright green fades away, we begin to see yellow and orange colors. Small amounts of these colors have been in the leaves all along - we didn’t them in the summer because they were covered up by the green chlorophyll. The bright reds and purples we see in leaves are made mostly in the fall. In some trees, like maples, glucose is trapped in the leaves after photosynthesis stops. Sunlight and the cool nights of autumn cause the leaves to turn this glucose into a red color. It’s the combination of all these things that makes the beautiful fall colors we enjoy each year.
Ever hear of Thomas Cole’s The Voyage of Life series? In 1840 he did this series of paintings that represent an allegory of the four stages, or seasons, of human life:
• In childhood, the infant glides from a dark cave into a rich, green landscape.
• As a youth, the boy takes control of the boat and aims for a shining castle in the sky.
• In manhood, the adult relies on prayer and religious faith to sustain him through rough waters and a threatening landscape.
• Finally, the man becomes old and the angel guides him to heaven across the waters of eternity.
In each painting, accompanied by a guardian angel, the voyager rides the boat on the River of Life. The landscape, corresponding to the seasons of the year, plays a major role in telling the story. And in those paintings you can clearly see the leaves changing colors in the season (manhood) that represents the fall of the voyager’s life.
So what’s this mean to you and me? Things change! Always! Life is full of changes and most of us are creatures of habit. And because we don’t know what’s next, we tend to cling to what we already have and know and are comfortable with. We reminisce about and cherish the past because it’s familiar, it’s already happened and we know how the movie ends. And while that’s generally true, it’s the half of the story that we tend to recognize. The other half is that the things we learn from the past should continually be updating our knowledge of life, and how to process the new things we see and experience, and how to better understand the meaning of who and what we are – that’s the harder part of the story to accept.
With each passing season, and the changes that occur, we need to grow and become wiser. And that wisdom should create the stuff we need to constantly be better, to do the things we’re called upon to do each day better, and to help those around us to become better. But you won’t learn anything or get better if you’re not open to the changes – natural or man-made – that occur every day.
I wish you could join me here at our camp to look across the lake at the beauty that is unfolding. The scene is constant; the colors let me know that time is marching on. On the one hand I could worry that the seasons of my life are marching on, or, on the other, I could be challenged by the things I’ve learned this year that will help me to be wiser and more thoughtful in the future. One stunts natural growth; the other invigorates a sense of wonder about the world around us and the endless possibilities that potentially exist. The choice is ours. And while these leaves will begin to fade and fall soon, the inspiration that they trigger should last a lifetime. That’s the voyage of life, and I’m sure glad to be on it!
My message this week is about being inspired to dream about improving our lives:
“You are never too old to set another goal or to dream a new dream.” -C.S. Lewis
Clive Staples Lewis (1898 – 1963), commonly referred to as C. S. Lewis and known to his friends and family as "Jack", was a British novelist, academic, medievalist, literary critic, essayist, lay theologian and Christian apologist from Ireland.
Got any new dreams today? Not the ones you try to remember and think about when you wake, but the kind that have you excited to try something really new. Everyone can dream, but not everyone has the curiosity, energy, courage and stamina to try to attempt and achieve their dreams. Most want things to be smooth and easy, with no surprises or challenges that can potentially make you look silly. Fact is, without those challenges or knowing how to recover from looking silly you’ll never get to experience what it is to learn from trying something new. You can tell the ones who are into this – the twinkle in their eye, the bounce in their step, the way they carry themselves. If that’s you, and you’ll know if it is, then set another goal today, dream another dream today and make a pledge to be creative and innovative today. Go ahead – you’re never too old!
Friday, September 30, 2011
at 5:24 AM
Friday, September 23, 2011
“Everyone wants to be true to something, and we’re true to you” - that’s the marketing tagline for Jet Blue’s travel rewards program. I know because it kept scrolling across the little screen on the back of the seat in front of me when I recently flew across country. It’s okay in the context of what they’re trying to promote, but it also might apply to more than just loyalty programs. And it may be that because people naturally want to be ‘true blue’ to so many things, it becomes overused and almost trite. That’s too bad. Because being ‘true blue’ can be a good thing.
First: ever wonder where the term ‘true blue’ comes from?
• Loyal and unwavering in one's opinions or support for a cause or product.
• 'True blue' is supposed to derive from the blue cloth that was made at Coventry, England in the late middle-ages. The town's dyers had a reputation for producing material that didn't fade with washing, i.e. it remained 'fast' or 'true'. The phrase 'as true as Coventry blue' originated then and is still used (in Coventry at least).
• True Blue is an old naval/sailing term meaning honest and loyal to a unit or cause.
• And dictionaries say that true blue refers to “people of inflexible integrity or fidelity”.
And second: does ‘true blue’ really mean anything in this era of fast food and slick advertising?
There are lots of loyalty programs – hotels, airlines, slot clubs, retail stores, pop food brands, credit cards, clothing, wine, restaurants, movie theaters, travel sites, theme parks, computer games and countless more – and they all try to get you to stick with them by rewarding you in all kinds of ways: points, miles, free gifts, shows, food and on and on. But it seems a bit contrived, as if there’s some Oz-like character behind a curtain trying to entice you with these awards (read: bribes).
Imagine if this kind of thing were done with going to school or work, singing in a choir, participating in some community event, volunteering your time to some worthy cause, remaining friends or staying in a relationship… doesn’t seem as appropriate in those, does it? Think of someone or something you really like: do you really and truly like them or it, or do you need to be bribed with rewards to feel that way. Of course you don’t. So why do the airlines and hotels and all those other things we purchase have to bribe us like them?
But – there are companies out there that do understand what it takes to win your loyalty:
• Southwest Airlines was one of the first companies that made having fun and using common sense part of their strategy for success. Singing the safety jingle, devising a different boarding routine and setting the record for on-time departures set them apart and won over customers. They got it!
• Zappos doesn’t give you anything extra to make you want to come back – they believe that great service plus free shipping and returns will do that. Everyone said that nobody would buy shoes online – wrong. Zappos gets it!
• Apple wins and keeps their customer’s loyalty by incubating and introducing cool new ideas and products all the time. And they’re just about the biggest and most successful and most admired company on the planet. They get it!
But for every Southwest Airlines-type great experience there are hundreds of others that under perform and underwhelm. So they sign you up and hope that rewarding your loyalty overcomes the other things they do that destroys your loyalty. Seems to me they just don’t get it?
Jet Blue says they give you more leg room – that’s true if you pay extra for those few rows that have it. How come they just don’t make eye contact and smile more? How come they can’t get the bags to the conveyor in less than 30 minutes (which may not seem like much to them but after a cross country flight an extra 30 minutes is painful). How come they don’t get it? I want to join their loyalty program so I can get another trip with them like I want to have my teeth drilled. And then they spend so much time and energy trying to give you that free round trip ticket if you apply for their credit card – you know, the one that has annual fees and high interest rates. How come they don’t get it? Why can’t they just treat me like a loyal and valued customer, like someone they genuinely like and appreciate, like they’d like to be treated if they had to fly on someone else’s airline. Seems to me they just don’t get it.
Most of the good things in life are rooted in quality, trust and respect. People you work with and for, family that you live with and love, things you do for fun and relaxation, games you gladly play with others, friendships you’re lucky enough to have, clubs you join and actively participate in, activities you sign up for – they’re all based on the simple premise that things that are good are that way because they are genuinely good and fun and worthwhile. And that’s why you stick with them loyally.
But all these other kinds of loyalty programs are contrived. And yet we sign up for them like they’re free and worthwhile. They’re not free – we pay for the increased costs of these rewards. And they’re not worthwhile - we’re treated poorly by those who have the attitude that the cheap rewards they give are enough to overcome the thoughtless and robotic service they go through the motions of providing. Next time someone asks if I’ve signed up for their loyalty program I’m going to give them a tip: treat me nicely, treat me fairly, treat me respectfully, act like you really do care, thank me like you really mean it and treat me like you really do want me as a customer – and I’ll come back as often as I can or need to, willingly and freely. When are all these marketing geniuses going to wake up? When are they going to be ‘true blue’ to the Golden Rule?
My message this week is about how excellence can lead to greatness:
”If you want to achieve excellence, you can get there today. As of this second, quit doing less-than-excellent work.” -Thomas J. Watson
Thomas John Watson, Sr. (1874 – 1956) was president of International Business Machines (IBM) and oversaw that company's growth into a global force from 1914 to 1956. Watson developed IBM's distinctive management style and corporate culture, and turned the company into a highly-effective selling organization. He was called the world's greatest salesman.
Do you want to achieve excellence? Some people don’t – they’re content to work alongside others, doing just enough to get by and satisfy their basic needs, content to have a few toys, take life easy and not make waves. But is that what you want – would that be enough for you? If not, then you’ve got to decide right now to start going farther, looking to help others, caring more, trying harder, and being more of what you can be today. You’ve got to take it to the next level – in commitment, in energy, in enthusiasm, in being a role model, in paying closer attention to details, in always striving to do and be all that you’re capable of. As of this second, you’ve got to quit doing less-than-excellent work. That’s how YOU can achieve excellence - (note: the emphasis is on YOU)!
at 5:34 AM
Friday, September 16, 2011
Where were you on 9/11? For most of us the answers are permanently etched in our minds. Like the attack on Pearl Harbor and VE Day for our parents, or the moment John Kennedy was shot or Armstrong set foot on the moon for the baby boomers, 9/11 has become one of the iconic moments in time for all who were alive then.
I remember exactly where I was, what I was doing, who told me and how I felt the day Kennedy was killed; and like most people I was watching on our little black and white TV when Ruby shot Oswald the next day. I remember my teacher bringing me into the assembly hall to watch when Armstrong took “one small step for man, one giant leap for mankind”. There have been literally trillions of moments in my life, but these iconic ones stand out, frozen in time and in my mind. And then there was 9/11.
In these weekly blogs I try to write about things that catch my attention. These stories tend to take on meanings beyond the specific incidents I mention, meanings that relate to life’s larger issues and that can possibly teach us something. But this one goes way beyond any of the moments and incidents that caught my attention - 9/11 caught the attention of everyone on the planet. There aren’t many things that reach that level, things that stop time, that leave indelible memories about where we were and who we were with, that immediately bring back visceral feelings and emotions of a long ago but clearly remembered moment in time. 9/11 does all of those things and more.
My wife and I were in NYC: preparing to get on the George Washington Bridge to go into Manhattan when the first plane hit; coming to a complete stop on the road and in our lives; watching in fear and confusion as the second plane hit; staring in horror as first one and then the other building fell; hearing about the other plane crashes in Washington and Pennsylvania; staying glued to the radio and then the television while the world stood still.
We drove away from the City that day in fear and confusion – trying to get as far away as possible and to make sense of how and why this happened. As we drove we came upon a rise in the road where all the cars were stopped; people were standing beside their cars and looking back in the direction we came from, so we stopped too. In the distance there was smoke where the towers so recently stood; nobody was talking; everyone was crying. We eventually made it to our home in the Adirondack Mountains, safe and overwhelmed by the fear and confusion that enveloped the world as we knew it. I can see and feel that day now as if were yesterday. I guess that’s what an iconic moment is: something we remember – clearly and forever.
And now, in what seems like no time at all, ten years have passed and the memorial to those killed has been unveiled. The reading of the names this past Sunday stopped and stunned us all over again. The tolling of the bells in New York, Washington and Shanksville brought us back to that moment in time. The sight of the grieving families and friends as they touched and etched the names of their fathers, brothers, mothers, sisters, relatives and friends brought us together now as we were back then. The pettiness and partisanship that dominates the news was pushed aside for just a moment as we all stood in solemn and shared tribute to something that transcended all the comparatively meaningless stuff that normally seeks to grab our attention. As sad as the memories are, the togetherness helps us get through the memories now like it did when this terrible tragedy first happened. Why can’t we make that feeling last?
A man named Al DiLascia from Chicopee, Mass. wrote a letter to the editor of the New York Times this week that summed this up:
For one brief moment on September 11, 2011, time seemed to stand still. People sought family members and recognized the importance of family. Acts of charity were plentiful. There was an assessment of life and what is really important. Places of worship were full. People unashamedly prayed. For one brief moment...
Let’s try to remember – not just the events that make up these iconic moments, but what they really mean, and what’s really important. Don’t let a day pass that you don’t tell those you love how much you care and to show it in thoughtful and meaningful ways, to touch the people and things that are most important to you, to reach out and give to those in need, and to quietly count and give thanks for all the blessings that are in your life. Do whatever you have to do to make the meaning of your iconic moments last!
My message this week is about being loyal to the people and things that are important in your life:
“Loyalty is something you give regardless of what you get back, and in giving loyalty, you're getting more loyalty; and out of loyalty flow other great qualities.”
Colonel Charles Edward ("Chuck") Jones (1952 – 2001) was a United States Air Force officer, a computer programmer, and an astronaut in the USAF Manned Spaceflight Engineer Program. He was killed in the attacks of September 11, 2001 aboard American Airlines Flight 11, the first plane to hit the first World Trade Center building at 8:46am.
All of the great values we read and write about seem to be interconnected, and loyalty may be the one at the hub of them all. Think of the people and things you’re loyal to, and then note the other great qualities that come from that loyalty. Friendship, success, pride, humility, professionalism, integrity, team spirit and passion are a few that immediately come to mind. These are the qualities and values that you hope to find in others, and certainly they’re the ones to which you should always aspire. But to get loyalty you need to give it, and that means you must be true to your work and family and friends, forgiving in your nature, humble in your approach to others, sincere in your dealings with all, and understanding in the complex and competitive world that we live in. Look for ways to give loyalty today without attaching any strings for reciprocity. And don’t be surprised if you then start to get loyalty and all the other great qualities flowing back to you in return.
Stay well. And please say a prayer for these heroes and all the others in your life who’ve passed.
at 6:20 AM
Friday, September 9, 2011
Vacation homes in the Adirondacks are commonly referred to as camps – my family is fortunate to have one and, as you know from some of my previous blogs, we’ve spent a lot of time there this year. These are not to be confused with day and overnight camps that parents send their kids to. This is about the second kind of camp.
I went to an overnight camp as a kid and loved it, but that’s a story for another time. This tale begins at Camp Nazareth (that’s the name of the overnight camp at the end of our lake). Its run by the local Catholic Diocese which has had little success in recent years attracting enough kids. More often than not, this wonderful facility – it can hold up to 300 kids at any one time - is terribly under used. Fortunately, it seems that they’ve now discovered ways to attract alternate users like family reunions, corporate retreats and, just this past week, a high school crew team (Google “rowing sport” to learn more about this sport on Wikipedia). And that crew team caught our attention.
Our family’s camp (we call it “The Point”) is on the water and we can easily see when anyone is on the lake. While sitting on our dock one morning we were surprised to see this crew team go by. If you’ve never seen a crew team before, they operate in long narrow boats (like large kayaks) that are referred to as “sculls” – these are two to eight-person boats that are rowed by that many team members, each of whom operates one oar. In this case, there were two eight-person sculls (one with all men and the other all women) that were practicing. Mind you, this is not an everyday sight – there are a few motorboats and a lot of canoes and kayaks on our lake, so the sight of these two sculls was a bit of a surprise. Alongside these two sculls was a small motorboat in which sat the coach who had a megaphone and was giving instructions and commands. On the first day of what appeared to be one of their initial practice sessions, these two sculls were having what was obviously some beginner’s training. And here’s another key bit of information: the team has to row in very close order for the boat to move along smoothly. If any of the rowers is out of synch (even a little) the boat can very easily (and visibly) miss a beat. And if any of those misses are overly pronounced the boats can stop altogether or even capsize. So at the beginning of this training the coach definitely wanted to take it slow.
As the week progressed, however, the boats began to move more smoothly, and over time they got smoother and faster. And since the object of crew is to beat the competition, smooth and fast is definitely better. In order to get smoother and faster, the individual team members all have to practice at learning not only how to improve their own skills but also how to be in better synch with all the other members of their team. In crew, as in so many other aspects of life, both are critical (as in one without the other is not worth much).
As we watched this unfold before us, we started to reflect on how the basic lessons being learned out on the lake apply to just about everything we do in life (and here I need to confess that my wife realized this before I did). Being effective and functional at anything – playing with friends on the school yard, getting along as a family, working with colleagues, participating on a sports team, singing in a choir, building something with others, participating in community events – really is about learning how to improve your own skills while also performing in concert with others. Learning anything alone is one thing, learning it together and then interacting with others is a whole different thing. The key to life is learning both, because one without the other is really not worth much. And here was a live metaphor for this right on the lake in front of us – and just like that my whole professional life flashed before me as I watched this training unfold.
Each of these young athletes was working hard to learn how to be the best they could be, they and their team mates were learning how to interact with each other more effectively, the coaches were seeing the results of their hard work and practice, and those of us on the sidelines were rewarded by seeing how things can and should work when effective instructions, practice and coaching all come together. We don’t often get to see things so clearly, or watch how the rituals of cause and effect play out so clearly. Simply put: this was a real lesson about life. And, in part because of where we were, and also because of what we saw and then realized, we were again moved to exclaim “that’s the Point!
My message this week is about finding things you can be passionate about, because they define who and what you are.
“I know that I have found fulfillment. I have an object in life, a task ... a passion.”
Amantine Lucile Aurore Dupin, later Baroness Dudevant (1804 – 1876), best known by her pseudonym George Sand, was a French novelist and memoirist.
Have you found fulfillment? Not just a momentary or fleeting sense of accomplishment, but a lasting and on-going feeling that “this is it”. We all do lots of little and mostly disconnected things – chores, work, hobbies – and these achieve short-term goals or complete individual assignments. But every now and then one big thing comes along that is more about defining our style or purpose, and these make us who and what we are. Now it could be a car or a job – those certainly say a lot about you. But to find fulfillment – to know that something is really about the “you” that is truly you – that’s a real find. And that’s the kind of thing that passion is truly built upon. Something you love deeply, that you can’t stop thinking about, that you can’t wait to get up and do each day, and that you truly care more about than almost anything else. That’s the kind of passion that is truly a treasure – and that’s the kind of object in life that you want to be on the lookout for – today and every day. That’s the Point!
at 5:14 AM
Friday, September 2, 2011
Last week was something else – an earthquake and a hurricane and tornados and sunshine and hot and cold… I'm having trouble remembering where I am.
I grew up in upstate New York and experienced four distinct seasons each year – but there were no earthquakes or tornados. I later moved to Nevada for nearly a quarter century and experienced dry heat – but there were never any hurricanes or tornados. I then moved to the beaches of California where the sun shines 300+ days a year, the temperature rarely gets above 75 and earthquakes and wild fires are a nuisance – but there are no tornados or hurricanes. And now I’m back in New York (city and upstate) and just about everything but wild fires have hit here in the past 8 months. What’s going on?
I didn’t own a winter coat – and the record snow falls and cold last winter drove me to Land’s End with a singleness of purpose. I didn’t own boots or an umbrella, and the wet snow and rains taught me a lot about what it means to stay dry. I’m used to driving wherever I want to go and not having a car here to help navigate through the varying weather patterns has made me a fan of the Weather Channel. I never thought about the weather, never worried about what I’d wear or looked at the skies for clues to what’s coming, and now that the weather changes in the blink of an eye I am obsessed with meteorology.
But last week, depending where you were in the path of all this weather, meteorologists either got it right, mostly right, or wrong. Hey – they’re human so maybe we shouldn’t hold them to such a high standard as always being right. I mean, is anybody always right? Maybe we should take what they say and apply some old fashioned lore to this inexact science – such as:
Red sky at night, sailor's delight,
Red sky in the morning, sailors take warning.
When the wind is blowing in the North
No fisherman should set forth,
When the wind is blowing in the East,
'Tis not fit for man nor beast,
When the wind is blowing in the South
It brings the food over the fish's mouth,
When the wind is blowing in the West,
That is when the fishing's best!
When halo rings the moon or sun, rain's approaching on the run.
When windows won't open, and the salt clogs the shaker,
The weather will favor the umbrella maker!
No weather is ill, if the wind be still.
When sounds travel far and wide,
A stormy day will betide.
If clouds move against the wind, rain will follow.
A coming storm your shooting corns presage,
And aches will throb, your hollow tooth will rage.
I wouldn’t normally be thinking about these things, but all this crazy weather has me spooked. Is it global warming or just the fact that weather seems unpredictable? Were the winters way more intense when we were kids, or did it just seem that way because we were kids? Can weather really be predicted correctly all the time by these meteorologists, or should we take what they say with a “grain of salt”? Or should we rely more on our own common sense as aided by some of these old fashioned sayings?
Here in New York last week the mayor and the meteorologists got it wrong – but not by much. The winds blew and the rains fell and, though there was less flooding and damage than predicted here, they made damn sure we were prepared by scaring the daylights out of us with their dire warnings. Now some people are complaining because they scared us; but those same people complained when they didn’t scare us before last winter’s massive snow storm, or that they didn’t scare others enough before Katrina.
Fact is, lots of people are never happy, especially if they’re inconvenienced. But potentially saving lives is better than trying to apologize for not saving lives: isn’t that what ‘better safe than sorry’ is all about? Maybe we expect too much from the elected officials who we don’t really like or trust anyways (especially when they are inconveniencing us). I guess they’re damned if they do and damned if they don’t. I’ve even read some editorials about how this should make us either for or against big government. Come on, it was just a storm. And even though lots of people got flooded out, and there was lots of damage to homes and fields and trees and power lines, and lots of high water and wind, I’m relieved because it was less than predicted here on my street. I’m really sad for those to whom it was as much or more than predicted. And even though I don’t blame anyone, I sure as hell would like to know what all this crazy weather means, and whether a red sky at night really does mean a sailor’s delight?
My message this week is about loyalty, and whether we need to think about how loyal we are to others and how loyal we need to be to ourselves:
“Loyalty to petrified opinion never yet broke a chain or freed a human soul.” -Mark Twain
Mark Twain achieved great success as a writer and public speaker. His wit and satire earned praise from critics and peers, and he was a friend to presidents, artists, industrialists, and European royalty.
Loyalty can be both good and bad. People often remain loyal long after the reason for doing so has ended. If the reason you became loyal has petrified then you need to re-examine your motives and goals; you need to break free when the times demand it and it’s the right thing to do. Loyalty should be given to the best ideas, the highest principles, the most ethical leaders, the greatest challenges, and to the most extraordinary opportunities. But sometimes we remain loyal just because we are afraid to appear disloyal or we’re afraid to re-examine that loyalty. This conflict can be a Catch 22, or it can be a moment of re-commitment and rebirth. And just like a plant that’s been sitting for a long time, it’s a good idea to re-pot our beliefs to make sure that our roots continue to grow deeper and stronger. So look at your loyalties today and make sure they’re where they should be.
Stay warm, dry and well!
at 5:36 AM |
Art. XXIX.—The Older Gravels of North Canterbury.
[Read before the New Zealand Institute, at Christchurch, 4th—8th February, 1919; received by Editor, 24th February, 1919; issued separately, 16th July, 1919.]
|General Description of the District where the Beds are best developed||269|
|Description of Typical Localities||270|
|Grey River, East Branch||270|
|Grey River, West Branch||273|
|Okuku River and Mairaki Downs||274|
|Kowai River, North Branch||274|
|Kowai River, South Branch||274|
|Lower Warpara Gorge||278|
|Other Canterbury Localities||279|
|General Conclusions as to the Origin and Age of the Beds, and Relation to the Gravels of the Canterbury Plains||280|
Widely distributed along the base of the Southern Alps lies a series of unfossiliferous sedimentary beds, consisting for the most part of well-stratified gravels, sands, and clays, with occasional lignite, whose position has hitherto been somewhat doubtful. Haast (1879, p. 316, and map, p. 370) included them in his Pareora formation, and mentioned the occurrence of lignite-beds (p. 318) in the “Moeraki” Downs, at the mouth of the Waipara River, and in the Broken River basin, but hardly mentions the locality where they attain their maximum development—viz., the Mount Grey Downs and the vicinity of the two branches of the Kowai River. Hutton (1885, p. 211) considered them as equivalent to the Wanganui system of the North Island, but remarked that they were difficult to distinguish from the upper gravels of the Pareora system. Park (1910, p. 252) considered them older fluvio-glacial drifts. Thomson (1917, p. 411) refers to them more fully, but is extremely doubtful whether they shall be assigned to his Notocene or Notopleistocene set of deposits.
Owing to the practical absence of fossils it is difficult to determine their position accurately, but they nevertheless represent an interesting series, and the following account is intended to bring out their chief features. In addition to the difficulty noted by Hutton, there is the additional one that in their lithological content as well as to some extent the conditions under which they were laid down they resemble the beds that overlie them, and this makes it at times impossible to separate them from subsequent gravel and sandy beds.
General Description of the District Where the Beds Are Best Developed.
The chief area where they are developed lies to the south-east and south of Mount Grey, between the Waipara and Okuku Rivers, but they attain their greatest development in the basins of the Southern Kowai and the Grey Rivers. Important outliers also occur to the west of the Okuku, on
the lower slopes of Mount Thomas, and south of the Ashley River, where they form the Mairaki Downs (= Moeraki Downs of Haast). The beds form a kind of frontal apron to the higher greywacke hills, such as Mount Grey and Mount Karetu; but still they rise in places to well over 1,000 ft. above sea-level. The downs country has been dissected to some extent, and on the front facing south-east consequent streams have cut deep narrow channels, with high precipitous banks, whereas in the north-eastern portion the tributaries forming the Northern Kowai tend to develop valleys along the strike. The same is also true of the east branch of the Grey. The character of the drainage points to recent and rapid uplift, perhaps in agreement with that of which there is distinct evidence on the coast farther north (McKay, 1877, p. 177; Hutton, 1877, p. 55; Speight, 1918, p. 99). On the sides of the steep banks, especially those running with the dip, numerous good sections are exposed; and it will be best at this stage to give a more detailed description of typical sections, preferably those illustrating the relationship between the underlying Tertiary beds and the overlying gravels. Although the beds are typically developed in the basin of the Kowai River, and I have selected the name of that locality as the one most appropriate to designate the series, yet the most instructive sections are to be seen in the basin of the Grey River, and these will therefore be taken first.
Descriptions of Typical Localities.
Grey River, East Branch.
The eastern or chief branch of the Grey River rises in the country between Mounts Grey and Karetu, flows south therefrom in a deep wooded gorge, and then gradually turns to the south-west and follows along the north-western edge of the Mount Grey Downs till it enters on the plains and joins the Okuku River in the neighbourhood of White Rock Station. The first part of its course has been cut in greywacke, but on leaving the higher country it crosses the marginal fringe of Tertiary sedimentaries at an angle of about 45° with their strike, so that when the stream runs in the direction of the dip the cross-section of its channel is narrow and trench-like, but when it runs along the strike the valley opens out somewhat, with dip slopes bordering the stream on its north-western side and steep scarp slopes on its south-eastern side. The latter are in places very bold and precipitous, and show clear-cut sections. Especially is this the case at the Horseshoe Cliff, about a mile below the gorge, where the north-western slope of the downs has been scored by a deep washout, and the strata are clearly exposed for 500 ft.
In the lower part of the river-gorge there is a most interesting occurrence of the lower members of the Cretaceo-Tertiary series, analogous to that seen in the Waipara and Weka Pass sections. These beds when followed along the strike run in the direction of the greywacke mass of Mount Grey; and unless they turn round on approaching it, as they do on the north-east slope of the mountain, the junction between the two sets of beds will in all probability be a fault contact. To the west of the gorge, however, the junction between the two sets of beds is a normal unconformity.
The following is a description of the beds here exposed, the sequence being in ascending order:—
Greyish sands and sandy shales, glauconitic, concretionary in places, and stained with sulphur; succeeded by light-coloured argillaceous and slightly glauconitic sands—all striking north-east and dipping south-east at 45°.
Greenish glauconitic sand passing up into glauconitic limestone, the glauconitic material being disposed in irregular patches and lenses, giving the rock a somewhat streaky appearance; it is also full of worm-borings filled with glauconitic material. This passes up into
Amuri limestone, 25 ft. thick, with less glauconite than 2. The passage beds between this and the lower bed consist of fragments of Amuri limestone in a greensand matrix, the limestone finally taking on the facies of the typical Amuri stone, being white and jointed into quadrangular blocks. The strike is as before, but the dip is less, being about 30°.
Glauconitic limestone, 20 ft. thick, comparable with the Weka Pass stone as it approaches a shore-line (Speight and Wild, 1918, p. 77), but passing up into a more sandy facies.
Marl, slightly sandy, with concretionary layers and rounded concretions. This is the stratigraphical equivalent of the “grey marl” in the Weka Pass district. It has the same strike and dip as the limestone, and its thickness is about 70 ft.
Thus far the sequence is quite clear and conformable, but for a time the exposures are obscured and the relations to the underlying beds are not plain.
Just below the gorge there is a well-marked bed, striking north-north-east, with slightly flatter dip than the limestone, and containing numerous specimens of Ostrea ingens. Farther down-stream, but higher in the series, is a sandy conglomerate followed by sands with broken shells. These pass up into sands with a layer of oyster and other shell fragments, and then follow the beds of the Kowai series.
These are first exposed at the mouth of the gorge, just above the site of the old sawmill. They consist of sands and sandy gravels containing shell-fragments and showing intraformational unconformities, but no clear evidence, given by sections, of an unconformity between the Kowai series and the lower Tertiaries. Almost everywhere in the case of gravels resting on sands or other finer detrital beds the upper surface of the latter has suffered some erosion, but in no case in this branch of the Grey does this, in my opinion, amount to sufficient to be considered a major unconformity.
On the next bluff down-stream, and higher in the sequence, the beds exposed consist of greenish-grey sands (weathering light-brown) and sandy gravels, with sandy carbonaceous shales and impure lignite. These are capped unconformably by terrace-gravels belonging to the early history of the Grey River. This sequence is repeated on the next bluff, but the gravel beds of the Kowai series become more important, one very heavy band of gravel near the top of the cliff being divided into two parts by a layer of carbonaceous shale. In the bed of the river, at the base of this cliff, is a section which shows an unconformable junction between a greenish sand and an overlying bed of gravel. After a careful consideration of the circumstances of this case I have come to the conclusion that it must be considered only as an intraformational unconformity, due to the erosion of the bed of sand by marine currents in the interval between its deposition and that of the succeeding layer of gravel.
These beds strike north-east, and dip south-east at an angle of 20°.
As the sections are followed down-stream their character does not change except that the gravels become increasingly important, a feature that is well exemplified at the Horseshoe Cliff, on the face of which gravels greatly predominate, some layers being from 50 ft. to 70 ft. in thickness. Well-defined sandy layers also occur. The regular stratification of the beds
towards the base of the cliff points to their having been deposited in shallow water in close proximity to a shore-line, and not on a land-surface; but at the higher levels the stratification becomes more indistinct and the pebbles become coarser and more subangular in shape, so it is almost certain that the closing beds of the series were laid down on a land-surface. The presence of lignite in the lower beds clearly indicates estuarine or deltaic conditions.
It should be noted that on the high banks of the Grey River there is a still more recent series of gravels belonging to the history of the stream. They are similar in lithological features to the gravels of the Kowai series, but they are neither so well stratified nor so well cemented. They are undoubtedly river and not sea deposits. Where contacts can be seen they are easily differentiated, but elsewhere, especially on the lower slopes of the downs, it is difficult to separate them from the upper members of the lower set of beds, which were also laid down on a land-surface.
Grey River, West Branch.
The general stratigraphy of the beds in the basin of the western branch of the Grey River is similar to that in the eastern. The following is a general description of the strata exposed above the greywacke as disclosed on the sides of the gorge of the stream:—
Sands and greensands.
Limestone, full of bryozoan remains, but only a few feet thick in the gorge of the stream, thickening, however, to the east and to the west. There is a marked difference in the features of this limestone as compared with that in the eastern branch, and as they are in apparent continuity it might be assumed that the stone in the western branch represents a shallower-water facies. I am by no means certain that this is the true explanation, and the question of the identity as regards their stratigraphical position must be reserved for further investigation.
Marls, greenish in colour, with rounded concretions and concretionary bands, passing up into greyish sands with fragmentary fossil shells.
In the river these beds strike E. 25° S. and dip south at an angle of 30°, but they have suffered some deformation, and the strike changes to north-east on the ridge between the two branches of the Grey, and also as the beds are traced round to White Rock and the Okuku River. The upper surface of the sands was distinctly eroded before the next bed was laid down. This consists of a heavy band of cemented gravel. The following beds are then encountered, in ascending order:—
Gravel bed just referred to.
Sandy clays and gravels.
Sandy clay and carbonaceous shale, repeatedly alternating. One bed of shale is from 12 in. to 18 in. thick.
Sandy gravel, well cemented with iron oxide.
Greenish-grey sands, sandy shales, and gravels, rapidly alternating, totalling over 200 ft. in thickness, the strike gradually becoming east-north-east, and the dip flattening out from 30° to 10°.
Gravels, sandy and with occasional thin layers of sandy clay, lying flat or with slight dip to the south-east. These are at least 500 ft. thick, and are well exposed on the ridge between the western Grey and the stream near the White Rock Station.
The section in this river thus shows that there is a distinct series in which gravels are the dominant beds lying unconformably on marine Tertiaries. It should be noted that in the western branch there are no gravel beds below the unconformity. Either they have never been deposited or they have been removed by erosion. There is a strong suggestion from the eastern branch that gravel beds are present among the higher members of the underlying marine series, so that their presence cannot be taken as decided evidence that beds containing gravels in this locality necessarily belong to the Kowai series.
Okuku River and Mairaki Downs.
Similar gravels occur on the banks of the Okuku, especially on the western side, where they form low hills fringing the base of Mount Thomas, and stretching westward towards the Garry River and Glentui. Towards the Ashley they are masked by more-recent gravels, but they reappear on the south bank of the river, forming the Mairaki Downs. The strata here consist of thick sandy gravels, sandy clays, and occasional layers of carbonaceous shale. Opposite the mouth of the Garry they strike north-east and dip north-west at an angle of 20°, forming the south-eastern wing of a syncline which is developed farther west, while farther east, towards Rangiora, the structure is anticlinal. The country directly between the Mairaki and Mount Grey Downs is probably a syncline, but the surface is completely masked by recent gravels and clays belonging to the Ashley and Okuku Rivers and to the lower course of the Grey and Makerikeri Rivers, the latter draining a considerable area on the south-western flank of the Mount Grey Downs.
Kowai River, North Branch.
An excellent idea of the structure and general features of the northern part of the downs area can be obtained by examination of the high banks of the North Kowai, and especially of a tributary which rises in Mount Brown itself and flows in a south-easterly direction across the strike of the beds, thus exposing all the members of the series present in this locality. The following is a general description of the beds encountered, starting with the Mount Brown beds and following up to the highest members of the series:—
At the contact with the upper members of the Mount Brown series the latter consist of sands, and marine gravels with shells, striking north-east and dipping south-east at an angle of 10°. The Mount Brown beds are here capped unconformably with sandy gravels containing rounded and sub-angular greywacke pebbles, and belonging in all probability to the high-level terrace-gravels of the present Kowai River. Lower down sands, sandy clays, and sandy gravels dipping south-east at very low angles are exposed on the banks of the stream and in the deep gullies on the northern side. There is certainly a disagreement in dip between these beds and the underlying Mount Brown beds, suggesting an unconformity, but nowhere could I see an actual contact in order to determine this point precisely. The slight escarpment of the downs which faces Mount Brown at this point is determined by the presence of the gravel beds which occur at this horizon. It is noteworthy that there is an entire absence of the gravel beds with broken-shell layers which cover the Mount Brown beds in the vicinity of Weka Pass, a point which increases the probability that the beds forming the downs rest on the Mount Brown beds unconformably.
Farther down-stream the beds lie almost flat, with an east-south-east strike and a dip to the north-north-east at very low angles (less than 5°). On a high bluff a series of well-stratified sands and sandy gravels is exposed. Near the base of the cliff, under a sandy bed cemented in its lower part with iron oxide, lies a narrow band of sandy carbonaceous shale, 6 in. to 8 in. thick, containing pieces of lignified wood, and passing down into sandy clay with interstratified irregular lenses of lignite. Under these lie sands and sandy gravels, and then bluish-green and brown sands. A little below this the strike swings round to north-north-east, with an easterly dip, and in a narrow gully on the south side of the stream an interesting section is exposed. Here both the bluish-green sands and the sandy gravels have been eroded, and on the eroded surface have been deposited sands and sandy gravels containing fragments of the lower beds. A similar occurrence is to be observed on the face of a cliff in the main stream, the lower beds dipping 10° and the upper lying flat across them. A thin layer of broken-shell fragments was observed high up on the face of the cliff in an inaccessible position.
Higher in the series are rapidly alternating sands and fine sandy gravels, in layers down to 1 in. in thickness, and these are succeeded by sands, sandy-gravel beds, and bluish-grey and brown sands, the former with broken-shell layers. In the gravels are numerous fragments of limestone, which must have been derived from a surface of the Amuri limestone exposed to decided erosion. The fragments are generally less than 2 in. in length, but are sometimes longer, and are usually flattened like beach shingle. There is no doubt as to the interstratification of these beds in the series under consideration, as the same feature was observed in a deep gully to the north of the stream in its proper stratigraphical position. The presence of these fragments is positive proof of the presence of an unconformity between these beds and the Amuri limestone, and supports the stratigraphical evidence from the Grey River. In the bluish-green sands there are occasional shell-fragments.
For some distance below this spot there are no clear sections, but sand is probably the major constituent of the beds. At the junction with the Kowai River, however, there are high cliffs on the northern side, where the strata are clearly visible for half a mile. The lowest beds exposed in this locality are sands with interstratified gravels, in which limestone-fragments form a most important constituent. The beds with the limestone-fragments are at least 50 ft. thick, and may be thicker. Higher up the limestone constituent gets less and less, and the pebbles are entirely of greywacke. No other included material, such as fragments of Mount Brown limestone, was noted at this spot, which might indicate the date of the break between the Amuri limestone beds and those under consideration. It is possible that these gravel beds are unconformable to the greenish marine sands, since for some distance no exposures are visible which enable their relations to be precisely determined, and there is evidence from other parts of the area that these upper gravels are unconformable to greenish sands—e.g., in the Grey River (see p. 272) and also in No. 2 Creek (see p. 276). These gravel beds are fairly well stratified, with occasional beds of sandy clay and thin carbonaceous shales, their total thickness being about 1,000 ft., and the whole thickness of the series from the junction of the Mount Brown beds upward being about 1,500 ft., though this may include two series—viz., the Motunau and the Kowai series.
The course of the main stream above its junction with the tributary follows almost along the strike of the beds, so that the structure is not so well displayed. The beds exposed consist of marine sands, which are remarkably
current-bedded, and loose and cemented sandy gravels with numerous fossils similar to the beds exposed in the Lower Waipara Gorge. The highest bed of this series exposed in the valley of the stream consists of greenish sandy clay, which weathers a light brown, and contains fossil shells. Its upper surface has been distinctly eroded, and on it rests a heavy layer of cemented gravel, and following this are sandy clays and gravel beds dipping south-east at angles of 10°. These beds pass upward into the gravel beds exposed on the cliffs of the river below the junction with the tributary. In the tributary mentioned above I could find no indication of an eroded surface analogous to that in the main stream, and so it may be an unconformity of local character similar to those recorded elsewhere, but it may indicate a decided unconformity between the Motunau and the Kowai series.
It is owing to the typical development of these gravels and the beds associated with them in the Kowai River, not only in this locality but in the south branch as well, that I have called them the Kowai series. It is possible, however, that the lower part of this group of beds may be equivalent to the upper part of the Motunau series, and subsequent investigation may show the term to be unnecessary.
In No. 2 Creek, a southern tributary of the North Kowai, there is a very important section. Just below the high bluff on the north side, about four miles above the junction, the stream has exposed the following beds:—
Greenish sands, becoming more clayey in the upper portions and passing up into sandy shale.
Lignite, very impure, 10 in. thick, striking east-north-east and dipping north-north-west at 5°.
Argillaceous sands, decidedly clayey above the coal but becoming more sandy and greenish in colour higher up. The thickness exposed is about 6 ft., but it is eroded, and sandy gravels rest on it unconformably. An eroded surface appears just below this in the bed of the stream, with an angular mass of green sandy clay embedded in the gravel.
In close proximity to the erosion surface there is another section showing the same features, but with only 3 ft. of bed 3 interposed between the coal and the gravels.
Just over the dividing-ridge between this and the South Kowai River there are high cliffs, facing south, composed of similar beds, with gravels more strongly developed in the higher levels, and dipping south at angles about 5°. Thus an anticlinal axis runs east-north-east along the ridge in close proximity to the road which runs along the crest.
Just at the point where the stream turns after leaving the steeper slopes of Mount Grey, and assumes a north-easterly course, coal and associated beds are exposed in its actual channel and in the bank of a small gully on the southern side. They consist of—
Greenish sands, passing up into sandy clays.
Clays succeeded by greenish sandy clays.
Gravels, mixed with sand, cemented with iron oxide.
Greenish sandy clays.
These beds strike north-east, and dip north-west at an angle of 5°.
When followed up-stream there is an alternation of sands and gravels, apparently conformable to the beds just enumerated, exposed in the slips
on the river-banks; but the dip becomes steeper till, on the face of a high bluff below the bush, it reaches 20°. Here are alternating sandy clays and gravels, the former greenish-yellow in colour, which are capped unconformably by somewhat irregular sands and gravels, lying almost horizontally across the denuded edges of the lower set. The upper series evidently forms the distinct ridge which leads down-stream past the point where the undoubted unconformity described above was observed.
Similar beds are observed in places on the banks of the stream higher up, but the covering of bush and soil is too complete to attempt a correlation with those lower down. Owing to this covering it is likewise impossible to say whether the junction between the greywackes of Mount Grey is a simple unconformity or a fault contact.
Kowai River, South Branch.
The high banks of this river rise in places to a height of 500 ft. above its bed, frequently with precipitous faces, and thus excellent sections are exposed. The strata are also folded into gentle anticlines and synclines, so that in the cores of the former the lower beds are exposed. They consist of the following in ascending order:—
Sands with concretionary layers, with broken-shell beds in the lower part, at least 80 ft. thick.
Green sandy clays and gravel beds, the latter finer in grain and thicker in the lower part, and cemented with iron oxide.
At higher levels there are rapidly alternating gravel and sandy beds, the former composed of subangular pebbles, which point to deposit either on a land-surface or on a shore-line in close proximity to the source of supply. No limestone pebbles were seen in these beds. In places they exhibit intraformational unconformities, such as one would expect when rapid changes in the conditions of deposit take place, especially when the change is from a sand to a gravel, and vice versa. The gravel beds frequently form steep cliffs; and their hard bands determine the dividing-ridges between the tributary streams running into the main river, especially on the south side, and they also determine an important reach of the river itself, although its direction is primarily across the strike, and therefore of consequent character.
A specially good section is to be seen where the river makes a right-angle turn, and changes from the subsequent to the consequent direction. The beds are here bent up into a rather sharp anticline with a north-east strike and a dip to the north-west at an angle of 50°. On the seaward side of the anticlinal axis the dip is much less, the angle being about 10°, the succession being similar. But farther down-stream the strike swings round till it is north-north-east, then north-north-west, and finally north-west, following for a time the direction of the main river. This allows the lower members of the series to be exposed again in the bed and banks of the river. They consist here of sands and sandy clays, greenish-blue in colour and weathering brown, containing fossils, some of the sands with concretionary layers and associated with thin gravel beds. The shells consist of Siphonalia, Glycymeris, and Ostrea, but in a fragmentary condition.
Farther down-stream the greenish-blue beds are still exposed, but the strike gradually becomes north-east with a dip to the south-east, and gravel beds form the greater part of the high bluffs which face the river on the north above the Mount Grey Station.
Opposite the right-angle turn of the river referred to above, and immediately to the west of the axis of the anticline, there is distinct evidence of the presence of an unconformity between the gravels just referred to and a higher series. The beds here consist of sands and sandy gravels which are lithologically indistinguishable from the higher members of the lower series. On a bluff facing the river the beds in contact with the gravels of the lower series are exposed lying across their denuded edges at low angles, and immediately to the north-west they show a reversal of dip and are inclined to the south-east at angles of from 5° to 7°. This dip is in agreement with that which can be observed in sections in the upper basin of the Southern Kowai, and notably to the south of the anticline which runs to the south-east of No. 2 Creek and parallel with it. Thus there is evidence of an unconformity in the Southern Kowai in close proximity to that in the No. 2 Creek in the drainage area of the Northern Kowai.
It will be noted that the lowest members of the series present in the two branches of the Kowai consist of sandy beds with marine fossils, whereas these do not appear with certainty in the Grey River. The gravel beds are, however, equally developed in each area. This difference is perhaps of no special stratigraphical importance, since the gravel beds in both areas are undoubtedly marine, and in the Grey area conditions may not have been favourable for the preservation of fossil remains. As far as I can see at present, there is no evidence of a major unconformity between the gravel beds and the lower marine beds, although minor, intraformational unconformities undoubtedly exist. There is, however, distinct evidence of a discordance at a higher level in the Kowai series between beds of similar lithological character both above and below the unconformity, but not such a discordance as necessitates the higher beds being placed in another distinct series. The whole area and its vicinity no doubt experienced a fairly rapid elevation, probably of a differential character, so that erosion went on in one part of the area while deposition was continuous in an adjacent part. If subsequent deposition over the denuded area then ensued this special stratigraphical feature can be satisfactorily explained.
In Fox's Creek, the next stream south of the Kowai, there is an excellent section of the gravels forming the great mass of the downs area. This stream rises in the centre of the area and flows east, being bounded on the north for the middle part of its course by precipitous banks, in places up to 500 ft. in height above the stream. The beds here exposed consist of sandy gravels, sands, and sandy clays, with occasional thin, discontinuous layers of sandy carbonaceous shale. The sands are frequently blue-green in colour, and without fossils as far as I could see. In this part of the course of the stream the beds lie flat, with a slight dip to the north, but on the eastern margin of the downs the dip increases to 10° and its direction becomes south-east. In the adjacent valleys to the north and south there are similar beds with similar dip.
Lower Waipara Gorge.
Just where the river crosses the western end of the Limestone Range, beds of the Motunau series are exposed, consisting of marine sands, sandy clays, and sandy gravels frequently cemented with calcareous material and containing numerous fossil shells (Speight, 1912 and 1914). They are
involved in anticlines and synclines, and sometimes dip at steep angles—as high as 55° to 60°. They do not contain, as far as my observation goes, any limestone-fragments such as might have been shed from the Weka Pass or Amuri limestone beds, and this suggests that they are not unconformable to the beds containing those limestones, a conclusion which is supported by general stratigraphical evidence. These Motunau beds are capped unconformably by gravels containing numerous fragments of limestone, and from their lack of distinct stratification it may be concluded that they are high-level terrace-gravels of more recent date. Similar gravels occur on the downs just east of the Amberley—Waipara Railway, covering a considerable extent of country, as they are occasionally exposed in the sides of deep washouts and cap the cliffs cut by the river in making its gorge.
Along the north bank of the river below the Teviotdale Bridge there is also a series of Motunau beds, consisting of sands, sandy clays, and gravels, some of which are very fine and smooth and are evidently of marine origin. These beds contain fossil shells at various levels, which point to the age being Mio-Pliocene or Pliocene (Speight, 1914). On the south bank of the river the beds are much obscured by slip-material and vegetation, so that in no place is the contact clearly displayed. On the terrace near the mouth of the river the following section is exposed:—
Yellowish sands, exposed at river-level and for 6 ft. upwards.
Sandy lignite, with well-marked woody structure and containing crystals of gypsum in stellate and columnar groups, 4 ft. in thickness.
Sandy clay (fireclay?), with occasional pieces of bituminized wood.
Lignite, full of bituminized wood, 6–8 in.
Sand and sandy clay, with pieces of wood, 4 ft.
Gravel, 6 ft.
Farther up-stream the lignite-beds are exposed in similar stratigraphical position, succeeded by yellowish clays and sandy gravels, and at one place there is an exposure under the lignite of well-stratified and rounded marine shingle.
These beds all strike east-south-east and dip west-south-west at very flat angles, and it is impossible to tell on stratigraphical grounds whether the gravels overlying the lignites are conformable or not, or whether they belong to the Kowai series or to recent terrace-gravels. The locality furnishes evidence of the ease with which conformity may be simulated under certain circumstances. If level beds are planed by the sea, and no irregularity left on their being depressed and covered with a veneer of sediments, apparent conformity may occur over considerable distances, and especially will this be the case if the beds in contact are of a sandy or gravelly nature. A suggestion of unconformity is given in this case by the presence in the lowest layer of gravel of large pebbles, up to 8 in. in diameter, and more or less subangular, indicating strong currents on a land-surface or on a sea-bottom in close proximity to land, and that the beds were deposited under conditions entirely dissimilar from those obtaining when the better-rounded gravels were laid down.
Other Canterbury Localities.
There are other localities in this part of the South Island where similar gravels occur, among which may be cited the Isolated Hills in the Culverden
Basin; the cliffs at Gore Bay, where well-cemented gravels are involved in a syncline; the western side of the Trelissick Basin, specially in the Hog's Back Creek; between the Pudding Stone and the North Ashburton River; and in South Canterbury between the Tengawai and Pareora Rivers. It is probable, too, that the deeper gravels encountered in the bore at Chertsey belong to this series, the indications of petroleum coming from plant-remains which elsewhere have formed lignite.* Just below the Rakaia Gorge, where great thickness of gravels has been exposed, there is an underlying set of beds which are more strongly oxidized than the covering strata, and they may perhaps be assigned to a series older than the prevailing shingle beds of the plains. Although this criterion is perhaps an unsatisfactory one on which to base a determination of relative age, yet it has been applied in Switzerland in order to differentiate the gravels of the older glacial series of that region.
General Conclusions as to the Origin and Age of the Beds, and Relation to the Gravels of the Canterbury Plains.
The materials of which these gravels are composed have been derived almost entirely from original greywackes. No limestone was noted among them except in the case of the beds in the North Kowai. Occasional pebbles of basalt also occur, such as might have been derived from areas where such rocks are known to exist. A siliceous sandstone, white in colour and forming rounded masses, which could not be traced to its source, also occurs freely in the gravels of the Mount Grey district. These are perhaps masses of sandstone which have been loosely cemented by processes analogous to those which have formed the sarsen stones, or “Chinamen,” as they are called by miners, of the schist areas of Central Otago.
The subangular nature of the pebbles shows that the greywacke land must have been in close proximity to the area of deposit. The absence of large pebbles suggests that it was of moderate relief, though it might have been the outlying portion of a more elevated tract. The gravels contain, however, no suggestion of a glacial or fluvio-glacial origin; they are just such gravels as might have been brought down by the present Ashley or Waipara Rivers, which have no connection with glaciers. The origin of the limestone constituent can be traced exactly, as exposures of limestone of similar nature occur within a short distance of the area where they have been deposited; but these limestone pebbles occur low down in the series and disappear at higher levels, so that the uppermost beds must have been derived from areas where limestone does not exist. Although the lower members of the series are undoubtedly marine, the upper members were in all probability deposited under estuarine conditions or actually on a land-surface.
The determination of the age of the Kowai series is a matter of some difficulty. The unconformity in the Grey River shows that it is certainly post-Miocene, and in the Lower Waipara Gorge beds occur under the gravels with a fossil content which shows them to be upper Pliocene (Speight, 1914, p. 300)—that is, beds which form the upper part of the Motunau series. Therefore we may reasonably infer that the Kowai series is either upper Pliocene, if no unconformity exists between two sets of beds, or Pleistocene, if an unconformity is demonstrable. The Pliocene beds of the lower Waipara are perhaps the uppermost beds of a conformable Cretaceo-Tertiary series.
[Footnote] * Pieces of carbonized bark have recently been obtained from a depth of 1,900 ft. in this bore.
Therefore if an unconformity can be proved between the Kowai series and any member of the lower series it will be unconformable to all. As the presence of included fragments of limestone in the gravel beds of the Kowai demonstrates the existence of a clear unconformity between the gravels and the limestones, and the undoubted erosion-surface in the Grey River demonstrates the presence of one at a higher level still, we may therefore infer that the gravels of the Kowai series must be of Pleistocene age. If, however, the Cretaceo-Tertiary series is broken up eventually into subordinate unconformable elements, then this argument fails, and the matter will depend on the relation of the Kowai series to the fossiliferous marine beds of the lower Waipara and the Northern Kowai, as being the highest beds on which the Kowai series undoubtedly rests. The relation of the two sets of beds is somewhat obscure, though, judging from the evidence in the latter locality, probably unconformable. Therefore all that can be definitely stated is that the Kowai series overlies undoubted upper Pliocene beds and must be of a later age, and is most probably Pleistocene.
This must be earlier than the gravels forming the Canterbury Plains, for these have suffered no deformation by folding movements, whereas the gravels of the Kowai series are at times folded somewhat acutely. They would therefore antedate the last period of glaciation to which the region had been subjected.
A point which bears on the conformity of the Tertiary sequence should be noted—viz., that in neither of the two branches of the Grey River are the typical Mount Brown beds developed, thus suggesting an unconformity between the Motunau series, or the Kowai series, and the Mount Brown beds in case the absence is due to erosion, or between the Mount Brown beds and the “grey marls” in case the overlying beds, together with the Mount Brown beds, are part of a conformable series. As there appears to be no evidence of unconformity between the Mount Brown beds and the “grey marls” in the typical locality, whereas there is some evidence of an unconformity between the Motunau series and the Mount Brown beds, it seems more likely that the absence of the Mount Brown facies in the Grey River is due to erosion of these beds after deposition. This, however, is a point which requires further investigation.
Haast, J. von, 1879. Geology of Canterbury and Westland.
Hutton, F. W., 1877. Rep. Geol. Explor. during 1873–74, pp. 27–58.
— 1885. Sketch of the Geology of New Zealand, Quart. Journ. Geol. Soc., vol. 41, pp. 191–220.
McKay, A., 1877. Rep. Geol. Explor. during 1874–76, pp. 172–84.
Park, J., 1910. The Geology of New Zealand, Whitcombe and Tombs, Christchurch.
Speight, R., 1912. A Preliminary Account of the Lower Waipara Gorge, Trans. N.Z. Inst., vol. 44, pp. 221–33.
— 1914. Additions to the List of Fossils from the Lower Waipara, Trans. N.Z. Inst., vol. 46, p. 300.
— 1918. Structural and Glacial Features of the Hurunui Valley, Trans. N.Z. Inst., vol. 50, pp. 93–105.
Speight, R., and Wild, L. J., 1918. The Stratigraphical Relationship of the Weka Pass Stone and the Amuri Limestone, Trans. N.Z. Inst., vol. 50, pp. 65–93.
Thomson, J. A., 1917. Diastrophic and other Considerations in Classification and Correlation, and the Existence of Minor Diastrophic Districts in the Notocene, Trans. N.Z. Inst., vol. 49, pp. 397–413. |
Part 2 - Those Who Are Unable to See the Fact of Creation
The theory of evolution is a philosophy and a conception of the world that produces false hypotheses, assumptions and imaginary scenarios in order to explain the existence and origin of life in terms of mere coincidences. the roots of this philosophy go back as far as antiquity and ancient Greece.
All atheist philosophies that deny creation, directly or indirectly embrace and defend the idea of evolution. the same condition today applies to all the ideologies and systems that are antagonistic to religion.
The evolutionary notion has been cloaked in a scientific disguise for the last century and a half in order to justify itself. Though put forward as a supposedly scientific theory during the mid-19th century, the theory, despite all the best efforts of its advocates, has not so far been verified by any scientific finding or experiment. Indeed, the "very science" on which the theory depends so greatly has demonstrated and continues to demonstrate repeatedly that the theory has no merit in reality.
Laboratory experiments and probabilistic calculations have definitely made it clear that the amino acids from which life arises cannot have been formed by chance. the cell, which supposedly emerged by chance under primitive and uncontrolled terrestrial conditions according to evolutionists, still cannot be synthesised even in the most sophisticated, high-tech laboratories of the 20th century. Not a single "transitional form", creatures which are supposed to show the gradual evolution of advanced organisms from more primitive ones as neo-Darwinist theory claims, has ever been found anywhere in the world despite the most diligent and prolonged search in the fossil record.
In their attempts to gather evidence for evolution, evolutionists have unwittingly proven by their own efforts that evolution cannot have happened at all!
The person who originally put forward the theory of evolution, essentially in the form that it is defended today, was an amateur English biologist by the name of Charles Robert Darwin. Darwin first published his ideas in a book entitled the Origin of Species by Means of Natural Selection in 1859. Darwin claimed in his book that all living beings had a common ancestor and that they evolved from one another by means of natural selection. Those that best adapted to the habitat transferred their traits to subsequent generations, and by accumulating over great epochs, these advantageous qualities transformed individuals into totally different species from their ancestors. the human being was thus the most developed product of the mechanism of natural selection. in short, the origin of one species was another species.
Darwin's fanciful ideas were seized upon and promoted by certain ideological and political circles and the theory became very popular. the main reason was that the level of knowledge of those days was not yet sufficient to reveal that Darwin's imaginary scenarios were false. When Darwin put forward his assumptions, the disciplines of genetics, microbiology, and biochemistry did not yet exist. If they had, Darwin might easily have recognised that his theory was totally unscientific and thus would not have attempted to advance such meaningless claims: the information determining species already exists in the genes and it is impossible for natural selection to produce new species by altering genes.
While the echoes of Darwin's book reverberated, an Austrian botanist by the name of Gregor Mendel discovered the laws of inheritance in 1865. Although little known before the end of the century, Mendel's discovery gained great importance in the early 1900s with the birth of the science of genetics. Some time later, the structures of genes and chromosomes were discovered. the discovery, in the 1950s, of the DNA molecule, which incorporates genetic information, threw the theory of evolution into a great crisis, because the origin of the immense amount of information in DNA could not possibly be explained by coincidental happenings.
Besides all these scientific developments, no transitional forms, which were supposed to show the gradual evolution of living organisms from primitive to advanced species, have ever been found despite years of search.
These developments ought to have resulted in Darwin's theory being banished to the dustbin of history. However, it was not, because certain circles insisted on revising, renewing, and elevating the theory to a scientific platform. These efforts gain meaning only if we realise that behind the theory lie ideological intentions rather than scientific concerns.
Nevertheless, some circles that believed in the necessity of upholding a theory that had reached an impasse soon set up a new model. the name of this new model was neo-Darwinism. According to this theory, species evolved as a result of mutations, minor changes in their genes, and the fittest ones survived through the mechanism of natural selection. When, however, it was proved that the mechanisms proposed by neo-Darwinism were invalid and minor changes were not sufficient for the formation of living beings, evolutionists went on to look for new models. They came up with a new claim called "punctuated equilibrium" that rests on no rational or scientific grounds. This model held that living beings suddenly evolved into another species without any transitional forms. in other words, species with no evolutionary "ancestors" suddenly appeared. This was a way of describing creation, though evolutionists would be loath to admit this. They tried to cover it up with incomprehensible scenarios. for instance, they said that the first bird in history could all of a sudden inexplicably have popped out of a reptile egg. the same theory also held that carnivorous land-dwelling animals could have turned into giant whales, having undergone a sudden and comprehensive transformation.
These claims, totally contradicting all the rules of genetics, biophysics, and biochemistry are as scientific as fairy-tales of frogs turning into princes! Nevertheless, being distressed by the crisis that the neo-Darwinist assertion was in, some evolutionist paleontologists embraced this theory, which has the distinction of being even more bizarre than neo-Darwinism itself.
The sole purpose of this model was to provide an explanation for the gaps in the fossil record that the neo-Darwinist model could not explain. However, it is hardly rational to attempt to explain the gap in the fossil record of the evolution of birds with a claim that "a bird popped all of a sudden out of a reptile egg", because, by the evolutionists' own admission, the evolution of a species to another species requires a great and advantageous change in genetic information. However, no mutation whatsoever improves the genetic information or adds new information to it. Mutations only derange genetic information. Thus, the "gross mutations" imagined by the punctuated equilibrium model, would only cause "gross", that is "great", reductions and impairments in the genetic information.
The theory of punctuated equilibrium was obviously merely a product of the imagination. Despite this evident truth, the advocates of evolution did not hesitate to honour this theory. the fact that the model of evolution proposed by Darwin could not be proved by the fossil record forced them to do so. Darwin claimed that species underwent a gradual change, which necessitated the existence of half-bird/half-reptile or half-fish/half-reptile freaks. However, not even one of these "transitional forms" was found despite the extensive studies of evolutionists and the hundreds of thousands of fossils that were unearthed.
Evolutionists seized upon the model of punctuated equilibrium with the hope of concealing this great fossil fiasco. As we have stated before, it was very evident that this theory is a fantasy, so it very soon consumed itself. the model of punctuated equilibrium was never put forward as a consistent model, but rather used as an escape in cases that plainly did not fit the model of gradual evolution. Since evolutionists today realise that complex organs such as eyes, wings, lungs, brain and others explicitly refute the model of gradual evolution, in these particular points they are compelled to take shelter in the fantastic interpretations of the model of punctuated equilibrium.
Is there any Fossil Record to Verify the Theory of Evolution?
The theory of evolution argues that the evolution of a species into another species takes place gradually, step-by-step over millions of years. the logical inference drawn from such a claim is that monstrous living organisms called "transitional forms" should have lived during these periods of transformation. Since evolutionists allege that all living things evolved from each other step-by-step, the number and variety of these transitional forms should have been in the millions.
If such creatures had really lived, then we should see their remains everywhere. in fact, if this thesis is correct, the number of intermediate transitional forms should be even greater than the number of animal species alive today and their fossilised remains should be abundant all over the world.
Since Darwin, evolutionists have been searching for fossils and the result has been for them a crushing disappointment. Nowhere in the world – neither on land nor in the depths of the sea – has any intermediate transitional form between any two species ever been uncovered.
Darwin himself was quite aware of the absence of such transitional forms. It was his greatest hope that they would be found in the future. Despite his hopefulness, he saw that the biggest stumbling block to his theory was the missing transitional forms. This is why, in his book the Origin of Species, he wrote:
Darwin was right to be worried. the problem bothered other evolutionists as well. A famous British paleontologist, Derek V. Ager, admits this embarrassing fact:
The gaps in the fossil record cannot be explained away by the wishful thinking that not enough fossils have yet been unearthed and that these missing fossils will one day be found. Another evolutionist paleontologist, T. Neville George, explains the reason:
Life Emerged on Earth Suddenly and in Complex Forms
When terrestrial strata and the fossil record are examined, it is seen that living organisms appeared simultaneously. the oldest stratum of the earth in which fossils of living creatures have been found is that of the "Cambrian", which has an estimated age of 530-520 million years.
Living creatures that are found in the strata belonging to the Cambrian period emerged in the fossil record all of a sudden without any pre-existing ancestors. the vast mosaic of living organisms, made up of such great numbers of complex creatures, emerged so suddenly that this miraculous event is referred to as the "Cambrian Explosion" in scientific literature.
Most of the organisms found in this stratum have highly advanced organs like eyes, or systems seen in organisms with a highly advanced organisation such as gills, circulatory systems, and so on. There is no sign in the fossil record to indicate that these organisms had any ancestors. Richard Monestarsky, the editor of Earth Sciences magazine, states about the sudden emergence of living species:
Not being able to find answers to the question of how earth came to overflow with thousands of different animal species, evolutionists posit an imaginary period of 20 million years before the Cambrian Period to explain how life originated and "the unknown happened". This period is called the "evolutionary gap". No evidence for it has ever been found and the concept is still conveniently nebulous and undefined even today.
In 1984, numerous complex invertebrates were unearthed in Chengjiang, set in the central Yunnan plateau in the high country of southwest China. Among them were trilobites, now extinct, but no less complex in structure than any modern invertebrate.
The Swedish evolutionist paleontologist, Stefan Bengston, explains the situation as follows:
The sudden appearance of these complex living beings with no predecessors is no less baffling (and embarrassing) for evolutionists today than it was for Darwin 135 years ago. in nearly a century and a half, they have advanced not one step beyond the point that stymied Darwin.
As may be seen, the fossil record indicates that living things did not evolve from primitive to advanced forms, but instead emerged all of a sudden and in a perfect state. the absence of the transitional forms is not peculiar to the Cambrian period. Not a single transitional form verifying the alleged evolutionary "progression" of vertebrates – from fish to amphibians, reptiles, birds, and mammals – has ever been found. Every living species appears instantaneously and in its current form, perfect and complete, in the fossil record.
In other words, living beings did not come into existence through evolution. They were created.
Deceptions in Drawings
The fossil record is the principal source for those who seek evidence for the theory of evolution. When inspected carefully and without prejudice, the fossil record refutes the theory of evolution rather than supporting it. Nevertheless, misleading interpretations of fossils by evolutionists and their prejudiced representation to the public have given many people the impression that the fossil record indeed supports the theory of evolution.
The susceptibility of some findings in the fossil record to all kinds of interpretations is what best serves the evolutionists' purposes. the fossils unearthed are most of the time unsatisfactory for reliable identification. They usually consist of scattered, incomplete bone fragments. for this reason, it is very easy to distort the available data and to use it as desired. Not surprisingly, the reconstructions (drawings and models) made by evolutionists based on such fossil remains are prepared entirely speculatively in order to confirm evolutionary theses. Since people are readily affected by visual information, these imaginary reconstructed models are employed to convince them that the reconstructed creatures really existed in the past.
Evolutionist researchers draw human-like imaginary creatures, usually setting out from a single tooth, or a mandible fragment or a humerus, and present them to the public in a sensational manner as if they were links in human evolution. These drawings have played a great role in the establishment of the image of "primitive men" in the minds of many people.
These studies based on bone remains can only reveal very general characteristics of the creature concerned. the distinctive details are present in the soft tissues that quickly vanish with time. with the soft tissues speculatively interpreted, everything becomes possible within the boundaries of the imagination of the reconstruction's producer. Earnst A. Hooten from Harvard University explains the situation like this:
Studies Made to Fabricate False Fossils
Unable to find valid evidence in the fossil record for the theory of evolution, some evolutionists have ventured to manufacture their own. These efforts, which have even been included in encyclopaedias under the heading "evolution forgeries", are the most telling indication that the theory of evolution is an ideology and a philosophy that evolutionists are hard put to defend. Two of the most egregious and notorious of these forgeries are described below.
Charles Dawson, a well-known doctor and amateur paleoanthropologist, came forth with a claim that he had found a jawbone and a cranial fragment in a pit in the area of Piltdown, England, in 1912. Although the skull was human-like, the jawbone was distinctly simian. These specimens were christened the "Piltdown Man". Alleged to be 500 thousand years old, they were displayed as absolute proofs of human evolution. for more than 40 years, many scientific articles were written on the "Piltdown Man", many interpretations and drawings were made and the fossil was presented as crucial evidence of human evolution.
In 1949, scientists examined the fossil once more and concluded that the "fossil" was a deliberate forgery consisting of a human skull and the jawbone of an orang-utan.
Using the fluorine dating method, investigators discovered that the skull was only a few thousand years old. the teeth in the jawbone, which belonged to an orang-utan, had been artificially worn down and the "primitive" tools that had conveniently accompanied the fossils were crude forgeries that had been sharpened with steel implements. in the detailed analysis completed by Oakley, Weiner and Clark, they revealed this forgery to the public in 1953. the skull belonged to a 500-year-old man, and the mandibular bone belonged to a recently deceased ape! the teeth were thereafter specially arranged in an array and added to the jaw and the joints were filed in order to make them resemble that of a man. Then all these pieces were stained with potassium dichromate to give them a dated appearance. (These stains disappeared when dipped in acid.) Le Gros Clark, who was a member of the team that disclosed the forgery, could not hide his astonishment:
In 1922, Henry Fairfield Osborn, the director of the American Museum of Natural History, declared that he had found a molar tooth fossil in western Nebraska near Snake Brook belonging to the Pliocene period. This tooth allegedly bore the common characteristics of both man and ape. Deep scientific arguments began in which some interpreted this tooth to be that of Pithecanthropus erectus while others claimed it was closer to that of modern human beings. This fossil, which aroused extensive debate, was popularly named "Nebraska Man". It was also immediately given a "scientific name": "Hesperopithecus Haroldcooki".
Many authorities gave Osborn their support. Based on this single tooth, reconstructions of Nebraska Man's head and body were drawn. Moreover, Nebraska Man was even pictured with a whole family.
In 1927, other parts of the skeleton were also found. According to these newly discovered pieces, the tooth belonged neither to a man nor to an ape. It was realised that it belonged to an extinct species of wild American pig called Prosthennops.
Did Men and Apes Come from a Common Ancestor?
According to the claims of the theory of evolution, men and modern apes have common ancestors. These creatures evolved in time and some of them became the apes of today, while another group that followed another branch of evolution became the men of today.
Evolutionists call the so-called first common ancestors of men and apes "Australopithecus" which means "South African ape". Australopithecus, nothing but an old ape species that has become extinct, has various types. Some of them are robust, while others are small and slight.
Evolutionists classify the next stage of human evolution as "Homo", that is "man". According to the evolutionist claim, the living beings in the Homo series are more developed than Australopithecus, and not very much different from modern man. the modern man of our day, Homo sapiens, is said to have formed at the latest stage of the evolution of this species.
The fact of the matter is that the beings called Australopithecus in this imaginary scenario fabricated by evolutionists really are apes that became extinct, and the beings in the Homo series are members of various human races that lived in the past and then disappeared. Evolutionists arranged various ape and human fossils in an order from the smallest to the biggest in order to form a "human evolution" scheme. Research, however, has demonstrated that these fossils by no means imply an evolutionary process and some of these alleged ancestors of man were real apes whereas some of them were real humans.
Now, let us have a look at Australopithecus, which represents to evolutionists the first stage of the scheme of human evolution.
Australopithecus: Extinct Apes
Evolutionists claim that Australopithecus are the most primitive ancestors of modern men. These are an old species with a head and skull structure similar to that of modern apes, yet with a smaller cranial capacity. According to the claims of evolutionists, these creatures have a very important feature that authenticates them as the ancestors of men: bipedalism.
The movements of apes and men are completely different. Human beings are the only living creatures that move freely about on two feet. Some other animals do have a limited ability to move in this way, but those that do have bent skeletons.
According to evolutionists, these living beings called Australopithecus had the ability to walk in a bent rather than an upright posture like human beings. Even this limited bipedal stride was sufficient to encourage evolutionists to project onto these creatures that they were the ancestors of man.
However, the first evidence refuting the allegations of evolutionists that Australopithecus were bipedal came from evolutionists themselves. Detailed studies made on Australopithecus fossils forced even evolutionists to admit that these looked "too" ape-like. Having conducted detailed anatomical research on Australopithecus fossils in the mid-1970s, Charles E. Oxnard likened the skeletal structure of Australopithecus to that of modern orang-utans:
What really embarrassed evolutionists was the discovery that Australopithecus could not have walked on two feet and with a bent posture. It would have been physically very ineffective for Australopithecus, allegedly bipedal but with a bent stride, to move about in such a way because of the enormous energy demands it would have entailed. By means of computer simulations conducted in 1996, the English paleoanthropologist Robin Crompton also demonstrated that such a "compound" stride was impossible. Crompton reached the following conclusion: a living being can walk either upright or on all fours. A type of in-between stride cannot be sustained for long periods because of the extreme energy consumption. This means that Australopithecus could not have been both bipedal and have a bent walking posture.
Probably the most important study demonstrating that Australopithecus could not have been bipedal came in 1994 from the research anatomist Fred Spoor and his team in the Department of Human Anatomy and Cellular Biology at the University of Liverpool, England. This group conducted studies on the bipedalism of fossilised living beings. Their research investigated the involuntary balance mechanism found in the cochlea of the ear, and the findings showed conclusively that Australopithecus could not have been bipedal. This precluded any claims that Australopithecus was human-like.
The Homo Series: Real Human Beings
The next step in the imaginary human evolution is "Homo", that is, the human series. These living beings are humans who are no different from modern men, yet who have some racial differences. Seeking to exaggerate these differences, evolutionists represent these people not as a "race" of modern man but as a different "species". However, as we will soon see, the people in the Homo series are nothing but ordinary human racial types.
According to the fanciful scheme of evolutionists, the internal imaginary evolution of the Homo species is as follows: First Homo erectus, then Homo sapiens archaic and Neanderthal Man, later Cro-Magnon Man and finally modern man.
Despite the claims of evolutionists to the contrary, all the "species" we have enumerated above are nothing but genuine human beings. Let us first examine Homo erectus, who evolutionists refer to as the most primitive human species.
The most striking evidence showing that Homo erectus is not a "primitive" species is the fossil of "Turkana Boy", one of the oldest Homo erectus remains. It is estimated that the fossil was of a 12-year-old boy, who would have been 1.83 meters tall in his adolescence. the upright skeletal structure of the fossil is no different from that of modern man. Its tall and slender skeletal structure totally complies with that of the people living in tropical regions in our day. This fossil is one of the most important pieces of evidence that Homo erectus is simply another specimen of the modern human race. Evolutionist paleontologist Richard Leakey compares Homo erectus and modern man as follows:
Leakey means to say that the difference between Homo erectus and us is no more than the difference between Negroes and Eskimos. the cranial features of Homo erectus resulted from their manner of feeding, and genetic emigration and from their not assimilating with other human races for a lengthy period.
Another strong piece of evidence that Homo erectus is not a "primitive" species is that fossils of this species have been unearthed aged twenty-seven thousand years and even thirteen thousand years. According to an article published in Time – which is not a scientific periodical, but nevertheless had a sweeping effect on the world of science – Homo erectus fossils aged twenty-seven thousand years were found on the island of Java. in the Kow swamp in Australia, some thirteen thousand year-old fossils were found that bore Homo Sapiens-Homo Erectus characteristics. All these fossils demonstrate that Homo erectus continued living up to times very close to our day and were nothing but a human race that has since been buried in history.
Archaic Homo Sapiens and Neanderthal Man
Archaic Homo sapiens is the immediate forerunner of contemporary man in the imaginary evolutionary scheme. in fact, evolutionists do not have much to say about these men, as there are only minor differences between them and modern men. Some researchers even state that representatives of this race are still living today, and point to the Aborigines in Australia as an example. Like Homo sapiens, the Aborigines also have thick protruding eyebrows, an inward-inclined mandibular structure, and a slightly smaller cranial volume. Moreover, significant discoveries have been made hinting that such people lived in Hungary and in some villages in Italy until not very long ago.
Evolutionists point to human fossils unearthed in the Neander valley of Holland which have been named Neanderthal Man. Many contemporary researchers define Neanderthal Man as a sub-species of modern man and call it "Homo sapiens neandertalensis". It is definite that this race lived together with modern humans, at the same time and in the same areas. the findings testify that Neanderthals buried their dead, fashioned musical instruments, and had cultural affinities with the Homo sapiens sapiens living during the same period. Entirely modern skulls and skeletal structures of Neanderthal fossils are not open to any speculation. A prominent authority on the subject, Erik Trinkaus from New Mexico University writes:
In fact, Neanderthals even had some "evolutionary" advantages over modern men. the cranial capacity of Neanderthals was larger than that of the modern man and they were more robust and muscular than we are. Trinkaus adds: "One of the most characteristic features of the Neanderthals is the exaggerated massiveness of their trunk and limb bones. All of the preserved bones suggest a strength seldom attained by modern humans. Furthermore, not only is this robustness present among the adult males, as one might expect, but it is also evident in the adult females, adolescents, and even children."
To put it precisely, Neanderthals are a particular human race that assimilated with other races in time.
All of these factors show that the scenario of "human evolution" fabricated by evolutionists is a figment of their imaginations, and that men have always been men and apes always apes.
Can Life Result from Coincidences as Revolution Argues?
The theory of evolution holds that life started with a cell that formed by chance under primitive earth conditions. Let us therefore examine the composition of the cell with simple comparisons in order to show how irrational it is to ascribe the existence of the cell – a structure which still maintains its mystery in many respects, even at a time when we are about to set foot in the 21st century – to natural phenomena and coincidences.
With all its operational systems, systems of communication, transportation and management, a cell is no less complex than any city. It contains power stations producing the energy consumed by the cell, factories manufacturing the enzymes and hormones essential for life, a databank where all necessary information about all products to be produced is recorded, complex transportation systems and pipelines for carrying raw materials and products from one place to another, advanced laboratories and refineries for breaking down imported raw materials into their usable parts, and specialised cell membrane proteins for the control of incoming and outgoing materials. These constitute only a small part of this incredibly complex system.
Far from being formed under primitive earth conditions, the cell, which in its composition and mechanisms is so complex, cannot be synthesised in even the most sophisticated laboratories of our day. Even with the use of amino acids, the building blocks of the cell, it is not possible to produce so much as a single organelle of the cell, such as mitochondria or ribosome, much less a whole cell. the first cell claimed to have been produced by evolutionary coincidence is as much a figment of the imagination and a product of fantasy as the unicorn.
Proteins Challenge Coincidence
And it is not just the cell that cannot be produced: the formation, under natural conditions, of even a single protein of the thousands of complex protein molecules making up a cell is impossible.
Proteins are giant molecules consisting of amino acids arranged in a particular sequence in certain quantities and structures. These molecules constitute the building blocks of a living cell. the simplest is composed of 50 amino acids; but there are some proteins that are composed of thousands of amino acids. the absence, addition, or replacement of a single amino acid in the structure of a protein in living cells, each of which has a particular function, causes the protein to become a useless molecular heap. Incapable of demonstrating the "accidental formation" of amino acids, the theory of evolution founders on the point of the formation of proteins.
We can easily demonstrate, with simple probability calculations anybody can understand, that the functional structure of proteins can by no means come about by chance.
There are twenty different amino acids. If we consider that an average-sized protein molecule is composed of 288 amino acids, there are 10300 different combinations of acids. of all of these possible sequences, only "one" forms the desired protein molecule. the other amino-acid chains are either completely useless or else potentially harmful to living things. in other words, the probability of the coincidental formation of only one protein molecule cited above is "1 in 10300". the probability of this "1" occurring out of an "astronomical" number consisting of 1 followed by 300 zeros is for all practical purposes zero; it is impossible. Furthermore, a protein molecule of 288 amino acids is rather a modest one compared with some giant protein molecules consisting of thousands of amino acids. When we apply similar probability calculations to these giant protein molecules, we see that even the word "impossible" becomes inadequate.
If the coincidental formation of even one of these proteins is impossible, it is billions of times more impossible for approximately one million of those proteins to come together by chance in an organised fashion and make up a complete human cell. Moreover, a cell is not merely a collection of proteins. in addition to proteins, cells also include nucleic acids, carbohydrates, lipids, vitamins, and many other chemicals such as electrolytes, all of which are arranged harmoniously and with design in specific proportions, both in terms of structure and function. Each functions as a building block or component in various organelles.
As we have seen, evolution is unable to explain the formation of even a single protein out of the millions in the cell, let alone explain the cell.
Prof. Dr. Ali Demirsoy, one of the foremost authorities of evolutionist thought in Turkey, in his book Kalitim ve Evrim (Inheritance and Evolution), discusses the probability of the accidental formation of Cytochrome-C, one of the essential enzymes for life:
After these lines, Demirsoy admits that this probability, which he accepted just because it was "more appropriate to the goals of science", is unrealistic:
The correct sequence of proper amino acids is simply not enough for the formation of one of the protein molecules present in living things. Besides this, each of the twenty different types of amino acid present in the composition of proteins must be left-handed. Chemically, there are two different types of amino acids called "left-handed" and "right-handed". the difference between them is the mirror-symmetry between their three dimensional structures, which is similar to that of a person's right and left hands. Amino acids of either of these two types are found in equal numbers in nature and they can bond perfectly well with one another. Yet, research uncovers an astonishing fact: all proteins present in the structure of living things are made up of left-handed amino acids. Even a single right-handed amino acid attached to the structure of a protein renders it useless.
Let us for an instant suppose that life came into existence by chance as evolutionists claim. in this case, the right and left-handed amino acids that were generated by chance should be present in nature in roughly equal amounts. the question of how proteins can pick out only left-handed amino acids, and how not even a single right-handed amino acid becomes involved in the life process is something that still confounds evolutionists. in the Britannica Science Encyclopaedia, an ardent defender of evolution, the authors indicate that the amino acids of all living organisms on earth and the building blocks of complex polymers such as proteins have the same left-handed asymmetry. They add that this is tantamount to tossing a coin a million times and always getting heads. in the same encyclopaedia, they state that it is not possible to understand why molecules become left-handed or right-handed and that this choice is fascinatingly related to the source of life on earth.13
It is not enough for amino acids to be arranged in the correct numbers, sequences, and in the required three-dimensional structures. the formation of a protein also requires that amino acid molecules with more than one arm be linked to each other only through certain arms. Such a bond is called a "peptide bond". Amino acids can make different bonds with each other; but proteins comprise those and only those amino acids that join together by "peptide" bonds.
Research has shown that only 50 % of amino acids, combining at random, combine with a peptide bond and that the rest combine with different bonds that are not present in proteins. to function properly, each amino acid making up a protein must join with other amino acids with a peptide bond, as it has only to be chosen from among the left-handed ones. Unquestionably, there is no control mechanism to select and leave out the right-handed amino acids and personally make sure that each amino acid makes a peptide bond with the other.
Under these circumstances, the probabilities of an average protein molecule comprising five hundred amino acids arranging itself in the correct quantities and in sequence, in addition to the probabilities of all of the amino acids it contains being only left-handed and combining using only peptide bonds are as follows:
As you can see above, the probability of the formation of a protein molecule comprising five hundred amino acids is "1" divided by a number formed by placing 950 zeros after a 1, a number incomprehensible to the human mind. This is only a probability on paper. Practically, such a possibility has "0" chance of realisation. in mathematics, a probability smaller than 1 over 1050 is statistically considered to have a "0" probability of realisation.
While the improbability of the formation of a protein molecule made up of five hundred amino acids reaches such an extent, we can further proceed to push the limits of the mind to higher levels of improbability. in the "haemoglobin" molecule, a vital protein, there are five hundred and seventy-four amino acids, which is a much larger number than that of the amino acids making up the protein mentioned above. Now consider this: in only one out of the billions of red blood cells in your body, there are "280,000,000" (280 million) haemoglobin molecules. the supposed age of the earth is not sufficient to afford the formation of even a single protein, let alone a red blood cell, by the method of "trial and error". the conclusion from all this is that evolution falls into a terrible abyss of improbability right at the stage of the formation of a single protein.
Looking for Answers to the Generation of Life
Well aware of the terrible odds against the possibility of life forming by chance, evolutionists were unable to provide a rational explanation for their beliefs, so they set about looking for ways to demonstrate that the odds were not so unfavourable.
They designed a number of laboratory experiments to address the question of how life could generate itself from non-living matter. the best known and most respected of these experiments is the one known as the "Miller Experiment" or "Urey-Miller Experiment", which was conducted by the American researcher Stanley Miller in 1953.
With the purpose of proving that amino acids could have come into existence by accident, Miller created an atmosphere in his laboratory that he assumed would have existed on primordial earth (but which later proved to be unrealistic) and he set to work. the mixture he used for this primordial atmosphere was composed of ammonia, methane, hydrogen, and water vapour.
Miller knew that methane, ammonia, water vapour and hydrogen would not react with each other under natural conditions. He was aware that he had to inject energy into the mixture to start a reaction. He suggested that this energy could have come from lightning flashes in the primordial atmosphere and, relying on this supposition, he used an artificial electricity discharge in his experiments.
Miller boiled this gas mixture at 100 0C for a week, and, in addition, he introduced an electric current into the chamber. At the end of the week, Miller analysed the chemicals that had been formed in the chamber and observed that three of the twenty amino acids, which constitute the basic elements of proteins, had been synthesised.
This experiment aroused great excitement among evolutionists and they promoted it as an outstanding success. Encouraged by the thought that this experiment definitely verified their theory, evolutionists immediately produced new scenarios. Miller had supposedly proved that amino acids could form by themselves. Relying on this, they hurriedly hypothesised the following stages. According to their scenario, amino acids had later by accident united in the proper sequences to form proteins. Some of these accidentally formed proteins placed themselves in cell membrane-like structures, which "somehow" came into existence and formed a primitive cell. the cells united in time and formed living organisms. the greatest mainstay of the scenario was Miller's experiment.
However, Miller's experiment was nothing but make-believe, and has since been proven invalid in many respects.
The Invalidity of Miller's Experiment
Nearly half a century has passed since Miller conducted his experiment. Although it has been shown to be invalid in many respects, evolutionists still advance Miller and his results as absolute proof that life could have formed spontaneously from non-living matter. When we assess Miller's experiment critically, without the bias and subjectivity of evolutionist thinking, however, it is evident that the situation is not as rosy as evolutionists would have us think. Miller set for himself the goal of proving that amino acids could form by themselves in earth's primitive conditions. Some amino acids were produced, but the conduct of the experiment conflicts with his goal in many ways, as we shall now see.
F Miller isolated the amino acids from the environment as soon as they were formed, by using a mechanism called a "cold trap". Had he not done so, the conditions of the environment in which the amino acids formed would immediately have destroyed the molecules.
It is quite meaningless to suppose that some conscious mechanism of this sort was integral to earth's primordial conditions, which involved ultraviolet radiation, thunderbolts, various chemicals, and a high percentage of free oxygen. Without such a mechanism, any amino acid that did manage to form would immediately have been destroyed.
F the primordial atmospheric environment that Miller attempted to simulate in his experiment was not realistic. Nitrogen and carbon dioxide would have been constituents of the primordial atmosphere, but Miller disregarded this and used methane and ammonia instead.
Why? Why were evolutionists insistent on the point that the primitive atmosphere contained high amounts of methane (CH4), ammonia (NH3), and water vapour (H2O)? the answer is simple: without ammonia, it is impossible to synthesise an amino acid. Kevin McKean talks about this in an article published in Discover magazine:
After a long period of silence, Miller himself also confessed that the atmospheric environment he used in his experiment was not realistic.
F Another important point invalidating Miller's experiment is that there was enough oxygen to destroy all the amino acids in the atmosphere at the time when evolutionists thought that amino acids formed. This oxygen concentration would definitely have hindered the formation of amino acids. This situation completely negates Miller's experiment, in which he totally neglected oxygen. If he had used oxygen in the experiment, methane would have decomposed into carbon dioxide and water, and ammonia would have decomposed into nitrogen and water.
On the other hand, since no ozone layer yet existed, no organic molecule could possibly have lived on earth because it was entirely unprotected against intense ultraviolet rays.
F in addition to a few amino acids essential for life, Miller's experiment also produced many organic acids with characteristics that are quite detrimental to the structures and functions of living things. If he had not isolated the amino acids and had left them in the same environment with these chemicals, their destruction or transformation into different compounds through chemical reactions would have been unavoidable. Moreover, a large number of right-handed amino acids also formed. the existence of these amino acids alone refuted the theory, even within its own reasoning, because right-handed amino acids are unable to function in the composition of living organisms and render proteins useless when they are involved in their composition.
To conclude, the circumstances in which amino acids formed in Miller's experiment were not suitable for life forms to come into being. the medium in which they formed was an acidic mixture that destroyed and oxidised any useful molecules that might have been obtained.
Evolutionists themselves actually refute the theory of evolution, as they are often wont to do, by advancing this experiment as "proof". If the experiment proves anything, it is that amino acids can only be produced in a controlled laboratory environment where all the necessary conditions have been specifically and consciously designed. That is, the experiment shows that what brings life (even the "near-life" of amino acids) into being cannot be unconscious chance, but rather conscious will – in a word, Creation. This is why every stage of Creation is a sign proving to us the existence and might of Allah.
The Miraculous Molecule: DNA
The theory of evolution has been unable to provide a coherent explanation for the existence of the molecules that are the basis of the cell. Furthermore, developments in the science of genetics and the discovery of the nucleic acids (DNA and RNA) have produced brand-new problems for the theory of evolution.
In 1955, the work of two scientists on DNA, James Watson and Francis Crick, launched a new era in biology. Many scientists directed their attention to the science of genetics. Today, after years of research, scientists have, largely, mapped the structure of DNA.
Here, we need to give some very basic information on the structure and function of DNA:
The molecule called DNA, which exists in the nucleus of each of the 100 trillion cells in our body, contains the complete construction plan of the human body. Information regarding all the characteristics of a person, from the physical appearance to the structure of the inner organs, is recorded in DNA by means of a special coding system. the information in DNA is coded within the sequence of four special bases that make up this molecule. These bases are specified as A, T, G, and C according to the initial letters of their names. All the structural differences among people depend on the variations in the sequence of these bases. There are approximately 3.5 billion nucleotides, that is, 3.5 billion letters in a DNA molecule.
The DNA data pertaining to a particular organ or protein is included in special components called "genes". for instance, information about the eye exists in a series of special genes, whereas information about the heart exists in quite another series of genes. the cell produces proteins by using the information in all of these genes. Amino acids that constitute the structure of the protein are defined by the sequential arrangement of three nucleotides in the DNA.
At this point, an important detail deserves attention. An error in the sequence of nucleotides making up a gene renders the gene completely useless. When we consider that there are 200 thousand genes in the human body, it becomes more evident how impossible it is for the millions of nucleotides making up these genes to form by accident in the right sequence. An evolutionist biologist, Frank Salisbury, comments on this impossibility by saying:
The number 41000 is equivalent to 10600. We obtain this number by adding 600 zeros to 1. As 10 with 11 zeros indicates a trillion, a figure with 600 zeros is indeed a number that is difficult to grasp.
Evolutionist Prof. Ali Demirsoy was forced to make the following admission on this issue:
In addition to all these improbabilities, DNA can barely be involved in a reaction because of its double-chained spiral shape. This also makes it impossible to think that it can be the basis of life.
Moreover, while DNA can replicate only with the help of some enzymes that are actually proteins, the synthesis of these enzymes can be realised only by the information coded in DNA. As they both depend on each other, either they have to exist at the same time for replication, or one of them has had to be "created" before the other. American microbiologist Jacobson comments on the subject:
The quotation above was written two years after the disclosure of the structure of DNA by James Watson and Francis Crick. Despite all the developments in science, this problem remains unsolved for evolutionists. to sum up, the need for DNA in reproduction, the necessity of the presence of some proteins for reproduction, and the requirement to produce these proteins according to the information in the DNA entirely demolish evolutionist theses.
Two German scientists, Junker and Scherer, explained that the synthesis of each of the molecules required for chemical evolution, necessitates distinct conditions, and that the probability of the compounding of these materials having theoretically very different acquirement methods is zero:
In short, the theory of evolution is unable to prove any of the evolutionary stages that allegedly occur at the molecular level.
To summarise what we have said so far, neither amino acids nor their products, the proteins making up the cells of living beings, could ever be produced in any so-called "primitive atmosphere" environment. Moreover, factors such as the incredibly complex structure of proteins, their right-hand, left-hand features, and the difficulties in the formation of peptide bonds are just parts of the reason why they will never be produced in any future experiment either.
Even if we suppose for a moment that proteins somehow did form accidentally, that would still have no meaning, for proteins are nothing at all on their own: they cannot themselves reproduce. Protein synthesis is only possible with the information coded in DNA and RNA molecules. Without DNA and RNA, it is impossible for a protein to reproduce. the specific sequence of the twenty different amino acids encoded in DNA determines the structure of each protein in the body. However, as has been made abundantly clear by all those who have studied these molecules, it is impossible for DNA and RNA to form by chance.
The Fact of Creation
With the collapse of the theory of evolution in every field, prominent names in the discipline of microbiology today admit the fact of creation and have begun to defend the view that everything is created by a conscious Creator as part of an exalted creation. This is already a fact that people cannot disregard. Scientists who can approach their work with an open mind have developed a view called "intelligent design". Michael J. Behe, one of the foremost of these scientists, states that he accepts the absolute being of the Creator and describes the impasse of those who deny this fact:
The result of these cumulative efforts to investigate the cell – to investigate life at the molecular level – is a loud, clear, piercing cry of "design!" the result is so unambiguous and so significant that it must be ranked as one of the greatest achievements in the history of science. This triumph of science should evoke cries of "Eureka" from ten thousand throats.
But, no bottles have been uncorked, no hands clapped. Instead, a curious, embarrassed silence surrounds the stark complexity of the cell. When the subject comes up in public, feet start to shuffle, and breathing gets a bit laboured. in private people are a bit more relaxed; many explicitly admit the obvious but then stare at the ground, shake their heads, and let it go like that. Why does the scientific community not greedily embrace its startling discovery? Why is the observation of design handled with intellectual gloves? the dilemma is that while one side of the elephant is labelled intelligent design, the other side must be labelled God.19
Today, many people are not even aware that they are in a position of accepting a body of fallacy as truth in the name of science, instead of believing in Allah. Those who do not find the sentence "Allah created you from nothing" scientific enough can believe that the first living being came into being by thunderbolts striking a "primordial soup" billions of years ago.
As we have described elsewhere in this book, the balances in nature are so delicate and so numerous that it is entirely irrational to claim that they developed "by chance". No matter how much those who cannot set themselves free from this irrationality may strive, the signs of Allah in the heavens and the earth are completely obvious and they are undeniable.
Allah is the Creator of the heavens, the earth and all that is in between.
The signs of His being have encompassed the entire universe.
1. Charles Darwin, the Origin of Species: By Means of Natural Selection or the Preservation of Favoured Races in the Struggle for Life, London: Senate Press, 1995, p. 134.
2. Derek A. Ager. "The Nature of the Fossil Record." Proceedings of the British Geological Association, vol. 87, no. 2, (1976), p. 133.
3. T.N. George, "Fossils in Evolutionary Perspective", Science Progress, vol.48, (January 1960), p.1-3
4. Richard Monestarsky, Mysteries of the Orient, Discover, April 1993, p.40.
5. Stefan Bengston, Nature 345:765 (1990).
6. Earnest A. Hooton, Up From the Ape, New York: McMillan, 1931, p.332.
7. Stephen Jay Gould, Smith Woodward's Folly, New Scientist, 5 April, 1979, p. 44.
8. Charles E. Oxnard, the Place of Australopithecines in Human Evolution: Grounds for Doubt, Nature, No. 258, p. 389.
9. Richard Leakey, the Making of Mankind, London: Sphere Books, 1981, p. 116
10. Eric Trinkaus, Hard Times Among the Neanderthals, Natural History, No. 87, December 1978, p. 10, R.L. Holoway, "The Neanderthal Brain: What was Primitive?", American Journal of Physical Anthrophology Supplement, No. 12, 1991, p. 94
11. Ali Demirsoy, Kalitim ve Evrim (Inheritance and Evolution), Ankara: Meteksan Yayinlari 1984, p. 61
12. Ali Demirsoy, Kalitim ve Evrim (Inheritance and Evolution), Ankara: Meteksan Yayinlari 1984, p. 61
13. Fabbri Britannica Science Encyclopaedia, Vol. 2, No. 22, p. 519
14. Kevin McKean, Bilim ve Teknik, No. 189, p. 7
15. Frank B. Salisbury, "Doubts about the Modern Synthetic Theory of Evolution", American Biology Teacher, September 1971, p. 336.
16. Ali Demirsoy, Kalitim ve Evrim (Inheritance and Evolution), Ankara: Meteksan Publishing Co., 1984, p. 39.
17. Homer Jacobson, "Information, Reproduction and the Origin of Life", American Scientist, January, 1955, p.121.
18. Reinhard Junker & Siegfried Scherer, "Entstehung Gesiche Der Lebewesen", Weyel, 1986, p. 89.
19. Michael J. Behe, Darwin's Black Box, New York: Free Press, 1996, pp. 232-233. |
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E-Cat Test Validates Cold Fusion Despite Challenges
The test of the E-Cat (Energy Catalyzer) that took place on October 6, 2011 in Italy has validated Andrea Rossi's claim that the device produces excess energy via a novel Cold Fusion nuclear reaction. Despite its success, the test was flawed, and could have been done in a way that produced more spectacular results -- as if confirmation of cold fusion is not already stunning enough.
Andrea Rossi stands in front of his E-Cat
apparatus, October 6, 2011
Photo by Maurizio Melis of Radio24
by Hank Mills
Pure Energy Systems News
Andrea Rossi has made big claims for the past year, about his cold fusion "E-Cat" (Energy Catalyzer) technology. He has claimed that it produces vast amounts of energy via a safe and clean low energy nuclear reaction that consumes only tiny amounts of nickel and hydrogen. A series of tests had been performed earlier this year that seemed to confirm excess energy is produced by the systems tested. Some of the tests were particularly impressive, such as one that lasted eighteen hours, and was performed by Dr. Levi of the University of Bologna. Unfortunately, the tests were not planned out as well as they could have been and had flaws.
The most recent test that took place on October 6, 2011 in Bologna, Italy, was supposed to address many of the concerns about the previous tests, and be performed in a way that would put to rest many issues that had been discussed continually on the internet. Despite showing clear evidence of excess energy -- which is absolutely fantastic -- this most recent test failed to live up to its full potential. It was a big success in that it validated the claim the E-Cat produces excess energy via cold fusion, but it was not nearly as successful as it could have been. Or as successful as we, the outsiders looking in, would like for it to have been.
The Inventor's Mindset
One thing that should be stated is that inventors do not always think like the people who follow their inventions. They have their own mindset and way of looking at things. This should be obvious, because they are seeing *everything* from a different perspective. For example, when we think seeing the inside of an important component would be exciting and informative, they consider it a threat to their intellectual property. Or, for example, when we would like to see a test run for days, they are thinking that a few hours is long enough. In their mind, they know their technology works, and running it for hours, days, or weeks would be more of a chore to them than an exciting event.
In Rossi's case, he has worked with these reactors for many years. He has tested them time and time again. In fact, he has built hundreds of units (of different models), and has tested every one of them. He is aware of how the units operate and how they perform. Actually, for a period of many months to a year or more, he had an early model of E-Cat heating one of his offices in Italy. Satisfying the curiosity of internet "chatters" by operating a unit for an extended period of time -- beyond what he thinks is needed to prove the effect -- is just a waste of his time, according to his thinking. He could spend the time getting the one megawatt plant ready to launch.
Don't forget, Rossi is a busy person. In addition to finishing the one megawatt plant, he has a new partner company to find, a wife at home, and a life to live! We need to consider that he works sixteen to eighteen hours a day building units, testing them, addressing other issues about the E-Cat. Although he is a very helpful person in many ways (willing to communicate with people and answer questions), he simply does not have the time to grant all of the many requests made of him. If he did, he could not get any work done at all, and the E-Cat would never be launched, or ever make it to the market place!
The Outsider's Mindset
I consider myself an outsider. I have never built a cold fusion device, have never spent years working to develop a technology, and have never gone through the grueling process of trying to bring a product into the market place. Although I spend a lot of time researching various technologies on the internet, I don't work sixteen to eighteen hours a day. In addition, I have no vested interest in the success of any technology, other than simply wanting at least one to hit the market place, ASAP.
As an outsider, I do not think like Rossi thinks. I don't think the majority of people think like Rossi thinks, because they are not in his shoes. They are not working to the point of exhaustion, and do not have years of their life invested in an exotic technology. Due to the fact we do not think like Rossi, his actions or sometimes lack thereof can seem strange, bizarre, or odd. Sometimes, they can make us want to smack ourselves, to make sure we are not in some sort of strange dream.
The recent test on October 6, 2011 is an example of a situation in which outsiders would have liked to have seen a very different test. Here are examples of how an outsider would have liked to see the test performed, compared to Rossi's possible mindset.
(Please note that I am making speculations about what Rossi is thinking, and his mindset. I do not know for sure if my guesses are accurate. If they are not, then I would like to apologize to Rossi, and give him the chance to respond in any way he sees fit.)
In the recent test, the output producing capability of the reactors was throttled down for safety reasons. This may have been done by keeping the hydrogen pressure low, or adding less of the catalyst to the nickel powder. Also, only one out of the three reactors that were inside of the module, were used in the test. For an experimental test to prove the effect beyond a shadow of a doubt, I, as an outsider, would have loved to have seen the device fully throttled up, despite the safety risks. Even if it meant everyone that attended would have had to sign long legal disclaimers, it would have been worth it.
I think it would have been great if all three reactors were utilized, and they all were adjusted to produce their maximum level of output. This would have increased the amount of output produced dramatically, and would have reduced the amount of input needed. The more heat produced by the system, the less heat would have needed to be input via the electric resistors.
Rossi, on the other hand, probably thought throttling up the device to a high level was not worth the risk, and was not needed to prove that excess energy was being produced. It is true that an explosion causing injuries -- while probably VERY unlikely -- could result in a setback of his project, and possible legal ramifications. Also, in reality, the test proved excess energy was being produced even with only one, throttled down reactor being used.
So even though a test of the device adjusted to operate at full power would have been useful and exciting, it was not absolutely needed for what Rossi wanted to accomplish.
I would like to ask Rossi to consider performing a demonstration with a module both adjusted to operate at full power, and utilizing all three reactor cores. Even if he has to limit the number of people involved, perform the test remotely with cameras monitoring the module, utilize a blast shield, or only allow certain individuals (who have signed disclaimers) to go into the room in which the module is running.
A Longer Self Sustain
As an outsider, I have not had the chance to look at test data from these devices self-sustaining for long periods of time -- 12 hours, 24 hours, days, weeks, etc. I would really like to see one of these units self sustain for a *very* long period of time. This is not because I think the output of the E-Cat during the recent test was due to stored energy being released (the 'thermal inertia' theory being floated around the internet). In fact, I think that the flat line in NyTeknik's graph -- showing self sustain mode for three and a half hours without any drop in output temperature -- provides clear evidence against the thermal inertia theory. The reason I would like to see a longer period of self sustain, is that it would not only document a huge gain of energy, but one that no individual could rationally deny!
Rossi has claimed that these devices represent an alternative energy solution that could change the world. I think this is true. However, to show just how much potential this technology has, an even more extended test of the E-Cat in self sustain mode (at full power or at least with all three reactors inside the module being used), would have been much more impressive. I am not saying the Oct. 6 test was not impressive -- it was very significant because it demonstrated excess energy and proof of cold fusion -- but that a longer test would have been better. It would have done more to shut up the cynics (a few of which will never change their minds), and help the technology get into the mainstream (dumbstream) media.
I really don't think Rossi cares too much about showing off the technology's full potential, at this point. He also does not seem to appear to want the attention of the mainstream media, or at least any more than he thinks he needs. If he did, the test would have been far different, and would have produced such a gigantic amount of excess energy everyone's jaws would have dropped. My jaw dropped when I saw the flat line during self sustain mode (because it proved beyond a doubt the system was producing excess energy), but my jaw did not drop as far as it could have, if the period of operation had lasted longer.
Interestingly, I have known inventors, of unrelated energy technologies, that purposely held back from showing the *best* version of their technology. They did not want to show off too much, because they did not want to deal with the fallout of attracting too much attention. Instead of performing an amazing demonstration, they performed one that proved the point -- at least to their satisfaction -- but would not attraction too much attention. I think Rossi may feel the same way. If he had his way, he would have never done a single test before the launch of the one megawatt plant. It was Focardi that convinced him to do a public test, because he feared that (due to health problems) he may not live long enough to see the technology be revealed to the world.
A longer test (at least 12 hours) in self sustain mode would have been great, exciting, and would have produced even more excess energy. However, in Rossi's mind, it was not needed, for potentially valid reasons (at least from the perspective of someone on the inside).
I would simply like to humbly plead with Rossi, to try and step in the shoes of the outsider, and at the next test allow the module to run for a longer period in self sustain mode.
Modern Testing Methods and Tools
I have looked at the data acquired during the test, but have not had a chance to study it as in depth as I would like to. The data shows a clear gain of energy in my opinion, and confirms that the E-Cat is producing excess energy. As I said before, the test was a success. However, it could have been performed in a more modern way.
For example, all of the temperature measurements, power input measurements, and water flow measurements should have been fed into the same computer, to be recorded in a real time manner. This way all the data would have been automatically recorded into one data set, including the hour and second of every measurement. It seems data collection was not done this way at the test, and some of the data was actually taken by hand!
Because the data was not all automatically recorded into one computer during the test, NyTeknik (who had the exclusive right to be the first to post a report on the test) has not yet posted a graph that charts all the measurements of all the factors of the test. What I would like to see is a single high resolution graph, that shows all of the measurements that were taken of every parameter of the test. If one graph showing everything would be too complex for a non-expert to easily interpret, then a series of graphs would be ideal. This would allow everyone to more simply determine the total energy in, and the total energy out.
The data collected and the manner in which it was collected is good enough to show there was a significant amount of excess energy produced, especially during self sustain mode. It may also be good enough to show even more details about the excess energy produced. Sadly, I'm not an expert in scientific data interpretation, so it takes me more time to interpret data than an expert who does so full time (like Rossi).
I hope that when I have had the time to examine the data in more depth, I will see that Rossi's claims about the results of the test (not just excess energy but a six fold gain of energy, in a worst case scenario) are accurate. At this point, I am not going to doubt him. He is the expert, and there are many people going over the data, and hopefully more data from the test will be coming in the near future.
What I would like to do, is request that he upgrade his data acquisition methods for any upcoming public tests. However, from Rossi's perspective, the way the data was acquired was good enough, and proved the point he wanted to make. I respect his view, but I do hope that he will change his mind in the future.
For the record, I am not stating that I think better data acquisition techniques are needed to verify his technology produces excess energy, and even significant amounts of it. I simply think it would make analysis of test data much simpler, quicker, and precise.
One of the most useful tools in the scientific method is a control. A control is an object or thing that you do not try to change during the experiment. For example, if you were giving an experimental drug to a hundred people, you might want to have a number of additional people who do not receive the drug. You would compare how the drug effects the people who consumed it, to those who did not receive the drug at all. By comparing the two sets of people, those who consumed the drug and those who did not, you could more easily see the effectiveness of the drug -- or if it was doing harm.
In Rossi's test, a control system would have been an E-Cat module that was setup in the exact same way, except it would have not been filled with hydrogen gas. It would have had the same flow of water going through it, the same electrical input, and it would have operated for the same length of time as the E-Cat unit with hydrogen. By comparing the two, you could easily see the difference between the "control" E-Cat (that was not having nuclear reactions take place), and the "real" E-Cat (that was producing excess heat).
If a control had been used in the experiment, the excess heat would be even more obvious. It would have been so obvious, that it could have made the test go from a major success (with some flaws), to the most spectacular scientific test in the last hundred years.
Yes, a control would have made that much of a difference!
I understand that Rossi may not see the need for a control, when the test that was performed clearly showed excess energy without it. A control might have made the experiment so mind blowingly amazing, it could have attracted too much media attention, too many scientists that would want to get involved, and too many individuals wanting additional information. The result could have been that Rossi would not even have the time to finish his one megawatt plant.
However, from the view point of an outsider, I think a control would have greatly benefited the experiment. If it created too much media attention, perhaps someone could volunteer to work for a month as an unpaid intern, filtering through all of the requests from media representatives, and taking care of many non-technical tasks, so Rossi could focus on getting the one megawatt plant ready!
I sincerely hope that during the test of the one megawatt plant, and any tests before then, a control run will be performed, in which no hydrogen is placed in the reactors.
Rossi's Statement about the Test Results
Andrea Rossi responded to an email we sent him that had questions about the test. Here is the email, and his responses.
THANK YOU FOR YOUR CONTINUOUS ATTENTION. PLEASE FIND THE ANSWERS IN BLOCK LETTERS ALONG YOUR TEXT:
Dear Andrea Rossi,
In regards to the latest test of the Energy Catalyzer, I have a number of questions I hope you can answer.
1) My understanding is that if a reactor core is not adjusted to be under-powered (below its maximum potential) in self-sustain mode, it can have a tendency to become unstable and climb in output. If the reactor is left in an unstable self-sustaining mode for too long, the output can climb to potentially dangerous levels. Can you provide some information about how the reactor core in the test was adjusted to self-sustain in a safe manner?
NO, VERY SORRY
a) For example, there was only one active reactor core in the module tested. How was the single reactor core adjusted to be under-powered?
b) Is adjusting the reactor core as simple as lowering the hydrogen pressure?
2) What is the power consumption of the device that "produces frequencies" that was mentioned in the NyTeknik article? Although the power consumption of this device is probably insignificant, providing a figure could help put to rest the idea (that some are suggesting) that a large amount of power was being consumed by the frequency-generating device, and transmitted into the reactor.
THE ENERGY CONSUMED FROM THE FREQUENCY GENERATOR IS 50 WH/H AND IT HAS BEEN CALCULATED, BECAUSE THIS APPARATUS WAS PLUGGED IN THE SAME LINE WHERE THE ENERGY-CONSUME MEASUREMENT HAS BEEN DONE
a) Can you tell us anything more about this frequency generating device and its function?
NO, SORRY, THIS IS A CONFIDENTIAL ISSUE
b) Is the frequency-generating device turned on at all times when a module is in operation, or only when a module is in self-sustain mode?
c) Some are suggesting that this device is "the" catalyst that drives the reactions in the reactor core. However, you have stated in the past that the catalyst is actually one or more physical elements (in addition to nickel and hydrogen) that are placed in the reactor core. Can you confirm that physical catalysts are used in the reactor?
YES, I CONFIRM THIS
3) Does the reaction have to be quenched with additional water flow though the reactor, or is reducing the hydrogen pressure enough to end the reactions on its own?
NEEDS ADDITIONAL QUENCHING
a) If reducing the hydrogen pressure (or venting it completely) is not enough to turn off the module, could it be due to the fact some hydrogen atoms are still bonded to nickel atoms, and undergoing nuclear reactions?
b) If there is some other reason why reducing hydrogen pressure is not enough to quickly turn off the module, could you please specify?
Thank you for taking the time to answer these questions, and for allowing a test to be performed that clearly shows anomalous and excess energy being produced. Hopefully, the world will notice the significance of this test.
THANK YOU VERY MUCH, AND, SINCE I HAVE ABSOLUTELY NOT TIME TO ANSWER (I MADE AN EXCEPTION FOR YOU) PLEASE EXPLAIN THAT BEFORE THE SELF SUSTAINING MODE THE REACTOR WAS ALREADY PRODUCING ENERGY MORE THAN IT CONSUMED, SO THAT THE ENERGY CONSUMED IS NOT LOST, BUT TURNED INTO ENERGY ITSELF, THEREFORE IS NOT PASSIVE. ANOTHER IMPORTANT INFORMATION: IF YOU LOOK CAREFULLY AT THE REPORT, YOU WILL SEE THAT THE SPOTS OF DRIVE WITH THE RESISTANCE HAVE A DURATION OF ABOUT 10 MINUTES, WHILE THE DURATION OF THE SELF SUSTAINING MODES IS PROGRESSIVELY LONGER, UNTIL IT ARRIVES TO BE UP TO HOURS. BESIDES, WE PRODUCED AT LEAST 4.3 kWh/h FOR ABOUT 6 HOURS AND CONSUMED AN AVERAGE OF 1.3 kWh/h FOR ABOUT 3 HOURS, SO THAT WE MADE IN TOTAL DURING THE TEST 25.8 kWh AND CONSUMED IN TOTAL DURING THE TEST 3.9 kWh. IN THE WORST POSSIBLE SCENARIO, WHICH MEANS NOT CONSIDERING THAT THE CONSUME IS MAINLY MADE DURING THE HEATING OF THE REACTOR DURING THE FIRST 2 HOURS, WE CAN CONSIDER THAT THE WORST POSSIBLE RATIO IS 25.8 : 3.9 AND THIS IS THE COP 6 WHICH WE ALWAYS SAID. OF COURSE, THE COP IS BETTER, BECAUSE, OBVIOUSLY, THE REACTOR, ONCE IN TEMPERATURE, NEEDS NOT TO BE HEATED AGAIN FROM ROOM TEMPERATURE TO OPERATIONAL TEMPERATURE.
WARMEST REGARDS TO ALL, ANDREA ROSSI
He claims that the test produced 25.8 kilowatt hours of power, and consumed only 3.9 kilowatt hours, not considering the losses from using two circuits of water and a heat exchanger. This would be very impressive for a system that is only using one reactor core (out of three), that has been adjusted to only produce a fraction of its maximum potential power.
However, from my analysis of the data so far (still trying to wrap my head around it), I have not been able to confirm his claim of a COP of 6. I am not saying it is not the case, or not in the data. I simply have yet to fully examine the data, and I am waiting for more data to be released.
Actually, I hope that someone will release all the data in one file and/or graph that will be easier to interpret. Perhaps NyTeknik, if they have not done so already, could contact Rossi or someone else who attended and recorded the data, and ask for any test data they are missing.
Bottom Line - Cold Fusion Is Here
The fact of the matter was that the October 6th test was a success in many ways.
- It documented a gain of energy.
- It documented a gain of energy in self-sustain mode.
- It documented massive "heat after death."
Most importantly, it proved beyond a doubt, that cold fusion is a reality.
Italian scientific journalist Maurizio Melis of Il Sole 24 Ore, who witnessed the test in Bologna,
"In the coming weeks Rossi aims to activate a 1MW plant, which is now almost ready, and we had the opportunity to inspect it during the demonstration of yesterday. If the plant starts up then it will be very difficult to affirm that it is a hoax. Instead, we will be projected suddenly into a new energetic era."
The test could have been made better in many ways. It had flaws. However, it was the most significant test of the E-Cat so far, for one reason in particular....
This graph shows that the E-Cat is a device producing excess energy, because the red line does not go down until after the hydrogen is vented.
- Some may legitimately argue about how much energy was produced, because we don't yet have all the test data in one easy to interpret graph or file.
- Some may point out the flaws in the test, such as the lack of a control, the lack of another several hours of operation in self sustain mode.
- Some may point out ways the test could be improved.
However, that graph by NyTeknik makes it clear the test was a success -- not a failure.
Mainstream media, your alarm clock is buzzing, it's time to wake up!
# # #
This story is also published at BeforeItsNews.
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Other PES Coverage
PESN Coverage of E-Cat
For a more exhaustive listing, see News:Rossi_Cold_Fusion
LENR-to-Market Weekly -- June 6, 2013 - EU Parliament gives LENR thumbs
LENR-to-Market Weekly -- May 30, 2013 - additional info on the E-Cat
3rd party test report (PESN)
LENR-to-Market Weekly -- May 23, 2013 - E-Cat 3rd Party Results Posted (PESN)
E-Cat Validation Creates More Questions (PESN;
May 21, 2013)
Third-Party E-Cat Test Results
Posted - posted on ArXiv.org
May 20, 2013)
Interview with E-Cat Distributor License Broker, Roger Green (PESN;
May 17, 2013)
LENR-to-Market Weekly -- May 9, 2013 - Interview with Rossi about recent 1 MW plant delivery
Interview with Andrea Rossi About 1 MW E-Cat Plant Delivery (PESN;
May 7, 2013)
LENR-to-Market MONTHLY -- April 29, 2013 - E-Cat teases with April 30
delivery date (PESN)
LENR-to-Market Weekly -- March 28, 2013 - E-Cat 3rd-Party testing
LENR-to-Market Weekly -- March 7, 2013 - more on NASA (PESN)
LENR-to-Market Weekly -- February 21, 2013 - NASA on nuclear reactor in
LENR-to-Market Weekly -- February 14, 2013 - Piantelli self-sustains 2
LENR-to-Market Weekly -- February 7, 2013 - CF 101 week 2 (PESN)
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By Isaac Sever, Cypress Semiconductor
Stepper motors convert electrical energy into discrete mechanical rotation. They are ideally suited for many measurement and control applications where positional accuracy is important. Stepping motors have the following advantages:
- Full torque when rotation is stopped. This is in contrast to brushed and brushless DC motors, which cannot provide full torque continuously when the rotor is stopped. This aids in maintaining the current position.
- Precise open-loop positioning and repetition. Stepper motors move in discrete steps as long as the motor stays under the maximum torque and current limits. This allows the rotor position to be determined by the control sequence without additional tracking or feedback. High quality stepping motors have three to five percent precision within a single step.
- Quick starts, stop, and reverse capability.
- High reliability because there is no brush or physical contact required for commutation. The life span of a stepping motor is dependent on the performance of the bearings.
- Microstepping mode can be used allowing direct connection to a load without intermediate gearing.
- A wide speed range can be controlled by varying the drive signal timing.
- Inherent resonance can cause noise, jerky rotation, and at extreme levels, loss of position.
- It is possible to lose position control in some situations, because no feedback is natively provided.
- Power consumption does not decrease to zero, even if load is absent.
- Stepping motors have low-power density and lower maximum speed compared to brushed and brushless DC motors. Typical loaded maximum operating speeds for stepper motors are around 1000 RPM.
- Complex electronic controls are required.
Figure 1: Structure of motors.
Types of stepping motors
There are several basic types of stepping motors:
- Variable reluctance motors with metal teeth.
- Permanent magnet motors.
- Hybrid motors with both permanent magnets and metal teeth.
Variable reluctance stepping motors have three to five windings and a common terminal connection, creating several phases on the stator. The rotor is toothed and made of metal, but is not permanently magnetized. A simplified variable reluctance stepping motor is shown in Figure 2. In this figure, the rotor has four teeth and the stator has three independent windings (six phases), creating 30 degree steps.
Figure 2: Simple variable reluctance stepping motor.
The rotation of a variable reluctance stepping motor is produced by energizing individual windings. When a winding is energized, current flows and magnetic poles are created, which attracts the metal teeth of the rotor. The rotor moves one step to align the offset teeth to the energized winding. At this position, the next adjacent windings can be energized to continue rotation to another step, or the current winding can remain energized to hold the motor at its current position. When the phases are turned on sequentially, the rotor rotates continuously.
The described rotation is identical to a typical BLDC motor. The fundamental difference between a stepper and BLDC motor is that the stepper is designed to operate continuously stalled without overheating or damage.
Rotation for a variable reluctance stepping motor with three windings and four rotor teeth is illustrated in Figure 3.
1, 2, 3, 1 → 3 steps → quarter turn
12 steps per rotation
As shown in Figure 3, energizing each of the windings in sequence moves the rotor a quarter turn, 12 steps are required for a full rotation.
Table 1: Variable reluctance stepper motor in Figure 3.
The three steps shown in Figure 3 move the rotor a quarter turn. A full rotation requires 12 steps for a variable reluctance stepper motor.
Typical variable reluctance motors have more teeth and use a tooth pole along with a toothed rotor to produce step angles near one degree.
Figure 3: Rotation control of variable reluctance stepping motor.
Permanent magnet stepping motor
A permanent magnet stepping motor consists of a stator with windings and a rotor with permanent magnet poles. Alternate rotor poles have rectilinear forms parallel to the motor axis. Stepping motors with magnetized rotors provide greater flux and torque than motors with variable reluctance. The motor, shown in Figure 4, has three rotor pole pairs and two independent stator windings, creating 30 degree steps.
Motors with permanent magnets are subjected to influence from the back-EMF of the rotor, which limits the maximum speed. Therefore, when high speeds are required, motors with variable reluctance are preferred over motors with permanent magnets.
Figure 4: Permanent magnet stepping motor.
Rotation of a permanent magnet stepping motor is produced by energizing individual windings in a positive or negative direction. When a winding is energized, a north and south pole are created, depending on the polarity of the current flowing. These generated poles attract the permanent poles of the rotor. The rotor moves one step to align the offset permanent poles to the corresponding energized windings. At this position, the next adjacent windings can be energized to continue rotation to another step, or the current winding can remain energized to hold the motor at its current position. When the phases are turned on sequentially the rotor is continuously rotated.
Rotation for a permanent magnet stepping motor with two windings and three pairs of permanent rotor poles (six poles) is shown in Figure 5.
Winding in sequence:
1 +/-, 2 +/-, 1 -/+, 2 -/+ → 3 steps → quarter turn
12 steps per rotation
Table 2: Permanent magnet stepping motor in Figure 5.
With one winding energized, the three steps move the rotor a quarter turn. A full rotation requires 12 steps for a permanent magnet stepper motor (bipolar) with both windings energized in each step. As shown in Figure 5, energizing each winding in sequence through each polarity moves the rotor a quarter turn. As before, 12 steps are required for a full rotation.
Figure 5: Rotation control of permanent magnet stepping motor, sequencing individual windings.
Another alternative to rotate a permanent magnet rotor is to energize both windings in each step. The vector torque generated by each of the coils is additive; this doubles the current flowing in the motor, and increases the torque. More complex control is also required to sequence the turning on and off of both windings.
As shown in Figure 6, energizing two windings in each step, sequencing through each combination of polarities moves the rotor a quarter turn. As before, 12 steps are required for a full rotation.
Table 3: Permanent magnet stepping motor in Figure 6.
Figure 6: Rotation control of permanent magnet stepping motor using both windings together.
Typical permanent magnet motors have more poles to create smaller steps. To make significantly smaller steps down to one degree, permanent magnet rotors can add metal teeth and toothed windings. This hybrid motor is described in the next section.
Hybrid stepping motor
Hybrid stepping motors combine a permanent magnet and a rotor with metal teeth to provide features of the variable reluctance and permanent magnet motors. Hybrid motors are more expensive than motors with permanent magnets, but they use smaller steps, have greater torque, and have greater maximum speeds.
A hybrid motor rotor has teeth placed on the directional axes. The rotor is divided into parts between constant magnet poles. The number of rotor pole pairs is equal to the number of teeth on one of the rotor’s parts. The hybrid motor stator has teeth creating more poles than just the main poles containing windings. The rotor teeth provide a smaller magnetic circuit resistance in some rotor positions, which improves static and dynamic torque. This is provided by corresponding teeth positioning; some parts of the rotor teeth are placed opposite the stator teeth and the remaining rotor teeth are placed between the stator teeth. Dependence between the number of rotor poles, the stator equivalent poles, and the phase number define step angle:
Figure 7: Hybrid stepping motor.
Rotation of a hybrid stepping motor is produced with the same control method as a permanent magnet stepping motor, by energizing individual windings in a positive or negative direction. When a winding is energized, a north and south pole are created, depending on the polarity of the current flowing. These generated poles attract the permanent poles of the rotor and the finer metal rotor teeth. The rotor moves one step to align the offset magnetized rotor teeth to the corresponding energized windings.
Stepping motor control
A step motor is a synchronous electric motor. Its fixed rotor equilibrium position occurs when aligned with the stator magnetic field. When the stator changes position, the rotor rotates to occupy a new equilibrium position.
There are several stepper motor drive modes:
- Full-step mode.
- Double-step mode.
- Half-step mode.
- Microstep Mode.
Stepping motors can be controlled in a variety of ways, trading off implementation requirements with greater accuracy and smoother transitions. Rotation control with full-steps, half-steps, and microsteps is described as follows:
Full-step mode for a permanent magnet and hybrid stepping motor is detailed in the Stepper Motor Introduction. Figure 5 illustrates one-phase full-step mode in which only one winding is turned on at a time. In this mode, the rotor’s balanced position for each step is in line with the stator poles. With only half of the motor coils used at a given time, the full torque obtained is limited.
Two-phase, full-step mode shown in Figure 6 uses both windings energized in each step. This doubles the current through the motor and provides 40 percent more torque than when only one phase is used at a time. With two windings energized, the rotor’s balanced position for each step is halfway between the two energized stator poles.
The full-step and double-step drive modes can be combined to generate half-steps of rotation for half-step mode. First one winding is turned on, and then the second winding is energized, moving the rotor half a step towards the second, as shown in Figure 8.
A half-step with the combination of one and two windings energized in full-step mode produces higher resolution, but does not provide constant torque throughout rotation.
Figure 8: Three half steps, 1/8 of a rotation.
Microstepping mode is an extension of the half-step drive mode. Instead of switching the current in a winding from on to off, the current is scaled up and down in smaller steps. When two phases are turned on and the current of each phase is not equal, the rotor position is determined by the current phase ratio. This changing current ratio creates discrete steps in the torque exerted on the rotor and results in smaller fractional steps of rotation between each full-step. Microstep mode reduces torque ripple and low-speed resonance present in the other modes and is required in many situations.
Microstepping creates rotation of the rotor by scaling the contributions of the two additive torque vectors from the two stepper motor windings.
Figure 9: Torque in microstepping control mode.
The total torque exerted on the rotor is the vector addition of the torque from the two rotors. Each of the torques is proportional to the position of the rotor and the sine/cosine of the step angle.
These equations can be combined and solved for the position of the rotor.
Fractional steps are created by scaling torque contributions between windings. Because torque is proportional to magnetic flux that is proportional to the current in the winding, the position of the rotor can be controlled by controlling the current flowing in each winding. To create smooth microsteps between full-steps, the current is varied sinusoidally with a 90 degree phase shift between the two windings as shown in Figure 10.
The current is scaled by controlling the root mean square (RMS) current using a current mode buck converter, commonly called a chopper drive when used with stepper motors. The phase current is converted into a voltage using a sense resistor in each phase ground path. This voltage is routed to a comparator that disables the output whenever the phase current rises above a reference. The comparator reference is provided by a voltage digital-to-analog converter (VDAC). By changing the VDAC supplied current limit for each microstep, the total motor torque remains approximately constant for each step of the sinusoidal current waveform.
Figure 10: VDAC current limit for microstep mode.
Microstepping allows the rotor position to be controlled with more accuracy and also has advantages in the rotation. Advantages of microstepping are:
- Position is controlled with more accuracy.
- Rotation can be stopped at specific fraction of a step.
- Transitions are smoother.
- Damp resonance creates fewer oscillations as motor steps (especially at startup and slowdown).
Figure 11: Smooth transitions between steps and limited oscillations and settling in microstep mode.
PSoC 3 introduction
The CY8C3866AXI device is in the PSoC 3 architecture. A block diagram of the device is shown in Figure 12 with the blocks used in the stepper application highlighted.
Figure 12: PSoC 3 (CY8C3866AXI) block diagram.
The PSoC 3 digital subsystem provides unique configurability of functions and interconnects. The stepper motor control uses these digital resources to implement timers, pulse width modulator (PWM) blocks, control registers, and a hardware lookup table (LUT).
The PSoC 3 analog subsystem provides the device the second half of its unique configurability. The stepper motor uses dedicated comparators, voltage DACs, and programmable gain amplifiers (PGA).
Stepper motor control based on PSoC 3
The block diagram of the stepper motor control based on the CY8C3866AXI is shown in Figure 13. The PSoC Creator™ schematic is shown in Figure 14.
Figure 13: Block diagram of PSoC 3 stepper motor controller.
Input control signals to the PSoC 3 device are:
- Motor Current Sensing: Analog input pins to detect motor phase current on shunt resistor. Used to limit current of the motor phases. See details in the following section.
- User Interface Pins
- User Input: Analog pin to read potentiometer for parameter input. Two digital pins for Menu control buttons.
- Character LCD: Digital output port (seven pins) to drive the character LCD on the DVK for menu options and user feedback.
- PWM signals to the high-side drivers (four digital output pins).
- PWM signals to the low-side drivers (four digital output pins).
Figure 14: PSoC Creator schematic for stepper motor control.
The PWMs are not used to produce the typical pulse width modulation output used with other motors. Instead, the PWMs act more as a timer to ensure a maximum chopping frequency to avoid overheating the drivers. Additionally, the PWM ‘kill circuit’ natively includes the cycle kill mode that implements the chopper drive method by disabling the drive outputs for the remainder of the current PWM cycle after the comparator trips.
The PWM signals are routed to a look up table (LUT) logic block, along with the current stepping stage index. This logic block implements a LUT using the PLD capabilities of a universal digital block (UDB) and routes the PWM signals to the eight legal output control combinations based on the current polarity of each phase. These control signals are routed through GPIOs to the external power driver circuits that drive the stepper motor. In the demonstrated chopper drive topology, transistors or MOSFETs are typically used to switch the high voltages and currents used to drive the stepper motors. The sequencing of the PWM control signals on the external power drivers produces the step by step rotation of the motor.
A timer generates periodic interrupts that generate each step (or microstep) of the motor. This timer can be used to run the motor at a specific speed, or to a specific position (exact number of steps). To set the speed of the motor, the interrupt period of the timer is updated by firmware.
PSoC 3 also implements current limiting for motor overcurrent protection and microstepping in hardware. This is described in the following section.
Microstepping and current protection implementation
Microstepping limits the current flowing in the motor windings to create smooth and well-controlled transitions between full-steps. This functionality also builds protection in hardware for overcurrent that shields the motor from damage. The block diagram of the system with the current feedback sensing paths is shown in Figure 15.
Figure 15: Overcurrent protection block diagram for microstepping.
Motor current is measured with two shunt resistors in the ground paths of the power driver MOSFETS (R1 and R2 in Figure 15). This voltage is low-pass filtered on the board and connected to two analog pins on PSoC 3 (labeled Curr_A and Curr_B).
The input voltages are fed into programmable gain amplifiers (PGA) implemented with the analog continuous time (CT) blocks. The PGA buffers the input voltage and drives it to a continuous time comparator. This voltage level from the sense resistor is compared to the current limit, set by an 8-bit voltage DAC. For microstepping the DACs’ output, sine and cosine waveforms are generated from a software lookup table. This limits the motor current sinusoidally for smooth microstepping.
The output of the comparator is connected to the PWM block and kills the PWM output when the current limit threshold is exceeded. This provides cycle-by-cycle current limiting to the motor and creates smooth microstepping transitions. The implementation of the current limiting protection in PSoC Creator is shown in Figure 16.
Figure 16: PSoC Creator schematic implementation of current limiting block for microstepping.
The PSoC 3 resources used in current limiting are:
- Two continuous time (CT) blocks implement the PGAs.
- Two fixed analog comparators are dedicated analog resources and do not use any SC/CT blocks.
- Two 8-bit PWMs implemented in UDBs (the same PWMs used to control the power device drivers). The output of the comparator triggers the kill input to the PWM when a current limiting condition is detected.
- Two 8-bit VDACs. These built in 8-bit voltage DACs are used to set the threshold for the comparator current limit.
Shown in Figure 17 are the settings for each DAC and rotation index, and the microstep pointer (ramping 1-128).
Figure 17: Current limit versus time step for 128 step microstepping.
The currents flowing in the two windings are measured with small sense resistors between the power devices and ground. The value of the current detection shunt resistor is a trade-off between power efficiency and robustness of the detection blocks. For a given current limit, enough change in voltage must be generated by the motor current to accurately detect the change with the comparator, but increasing the resistor increases heat and reduces efficiency.
The current limiting protection mechanism implemented in PSoC 3 hardware is an on-chip low-cost solution.
The output PWM drivers are controlled by a hardware lookup table. The table takes inputs from the two PWM blocks and a control register that holds the rotation index (as shown in Figure 18).
In Table 4, the PWM control hardware LUT receives the stage index and PWM signals as inputs and outputs the eight PWM driver signals.
Table 4: PWM control hardware LUT.
Figure 18: PSoC Creator schematic LUT implementation of PWM output control.
When operating with microstep drive mode, the PWM outputs for PWM_A and PWM_B cycle through 01, 10, and 11. When the stepper motor operates under full-step mode, both PWMs are on (11). In this case, the LUT simplifies to the following table the rotation sequence described in the full-step descriptions earlier.
In Table 5, the simplified MPhase output control hardware LUT receives the stage index and PWM signals as inputs and outputs the eight PWM driver signals.
Table 5: MPhase output control hardware LUT.
The stepper motor can be run at a fixed speed or to a desired position. To run at a fixed speed, the timer period that triggers each step (or microstep) is adjusted. The 16-bit timer terminal count triggers an interrupt that is used to initiate each step. Input frequency of the timer is 100 kHz to ensure precision speed control. PSoC 3 is also able to receive step pulse commands from an external controller such as a PLC.
In Figure 19, the timer terminal count triggers an interrupt that initiates each step.
Figure 19: PSoC Creator schematic implementation of speed control timer.
To run in position control mode, the stepper motor turns a specific number of steps and then stops. (Position control mode is not supported with the user interface in the stepper motor demo). An internal counter is used to count the desired steps. When the desired position is reached, the step control from the timer interrupt is masked until the user requests another action.
When the motor stops, phase current is lowered automatically to save power and reduce heating.
The ability to control the position in an open loop configuration by counting steps (or microsteps) is dependent on the stepper motor operating within the torque and motor load limits. If the torque/load limits are exceeded, the motor can miss steps and the absolute rotational position information is lost.
There is one main loop and one interrupt service routine (ISR) for control of the motor, the timer ISR. The timer ISR generates an interrupt that triggers the step control function (see Figure 19). Each time the step function is called, the motor takes one step (or microstep). The step function looks up the sinusoidal values from a table and sets the DAC output voltage to control the phase currents. A flow chart of the firmware operation is shown in Figure 20. Other ISRs for the UART and ADC are also used for the demo project UI and GUI interfaces.
Figure 20: Stepper motor control firmware flow chart.
PSoC resource utilization
The stepper motor uses resources from the digital and analog portions of the PSoC 3 device. The highest use of resources stem from the VDACs and comparators. Two VDACs and two comparators are used for the stepper motor microstepping control. This constraint limits the CY8C3866AXI-040 device to a maximum of two stepper motor controllers.
Table 6: Stepper motor demo CY8C3866AXI-040 resource utilization (blocks with none used are not shown).
Table 7: Stepper motor demo on CY8C3866AXI-040 memory utilization (Keil™ Complier, Level-5 optimization)
Cypress’ stepper motor control with PSoC 3 incorporates current limiting and microstepping control for an optimized solution. Up to 128 microsteps is suitable for precision position control. The PSoC 3 stepper motor control solution has low total system cost and leaves significant PSoC 3 resources available for additional system functions.
- Cypress Application Note AN2229, “Motor Control - Multi-Functional Stepping Motor Driver” by Victor Kremin and Ruslan Bachinsky.
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SINCE evidence of the existence of the belief that the soul lives on is so indubitable, the question arises--what is its condition? In what state does the discarnate spirit find itself after final separation from the body? And first, as to what we may be allowed to call, for want of a better term, its physical condition.
We have already noted that soul is conceived as having both form and substance, the latter, so to speak, greatly rarefied. Moreover, it has been brought to our attention that the most common idea concerning form is that the soul is a replica of the body it inhabited. Consistency in primitive thinking is not to be assumed, as we have seen, nor are logical processes among primitives quite the same as ours. Yet when a disembodied soul took up its post-mortem residence in a serpent, for example, we may not suppose that that soul was still regarded as human in shape. But so far as the author has discovered, no decisive evidence exists on this point. The probabilities favor greatly the supposition that in such cases transformation of the soul shape was supposed to have taken place. Evidence of the common idea, retention by the soul of its human shape, has been before us. We have noted that some tribes mutilate the body of the dead, thinking that by so doing they inflict like wounds upon the soul and thus impose incapacity for harm upon the ghost, the double of the body. The Omahas slit the soles of a murdered man's feet that his spirit may be unable to return and cause damage to the people. Mangaeans prefer death in battle--men are then in their full strength; disease weakens them, and souls have the nature of the body at death. Barongo believe that souls are young or old, according to the age at death, and so do the Indians of Gran Chaco. Naga tribes of Manipur think that ghosts bear whatever tattoo marks, mutilations, or other blemishes or embellishments occurred on the body. Some people carry this idea so far as to prefer
[1. Fletcher and La Flesche, 27th Report, etc., p. 215.]
death before decay of natural powers sets in, and so commit suicide or are buried alive, that the soul may continue to exist in full vigor.
Having form and substance, the soul has certain physical needs. It hungers, thirsts, feels cold and heat. The degrees of grossness of these wants vary greatly. Sometimes the hunger, thirst, and wants and passions may be appeased by the mere spirit or ghost of food, drink, etc.; and the ghosts are served by the spirits or (as our theosophical friends might be imagined as saying) the astral bodies of dishes, implements, or weapons which are destroyed (i.e., killed) that their spirits may accompany the ghost into the spirit land. Indeed, this is by all odds the most prevalent conception. Sometimes it is the more evanescent or the more vital elements, such as the blood, which are used by the ghost, as in the celebrated case of Tiresias in the Odyssey. The cases already cited of food, drink, weapons, utensils, and the like possessing souls and being offered or placed with the dead, oftentimes being broken or mutilated so as to "kill"
[2. Cases are cited in Frazer's Dying God, pp. 9-14.
3 Book XI.]
them, furnish direct testimony to the supposed needs of the ghost. The hunger felt by the disembodied soul is vividly expressed by most African tribes, whose belief is that ghosts can and do eat even human bodies A Ghosts also suffer from cold, hence New Guineans, and others, make fires at the graves, and even build huts, so that when the ghosts come up from the body they may find comfort.'
Ghosts have voices, too, but thin and shadowy like themselves. They chirp like crickets or utter their words in whistling tones. So the wizards by ventriloquistic art impose upon the credulous, and by wheezing utterance produce the effect of communications from a shadowy being or from the ground. Note the indications of shamanistic practice in the Prophet Isaiah (8:19- 29:4).
What we may regard as the disposition of the ghost is by most peoples held to be fixed by the character of the person while on earth.
[4. Talbot, In the Shadow of the Bush, pp. 224-225, 232-233, 238, etc.; ERE, vi. 65 ff. The testimony is being exhaustively collected in Frazer, Belief in Immortality--see the Index, under "Food."
5. Brown, Melanesians and Polynesians, pp. 442 ff.; Neuhass, Deutsch Neu-Guinea, iii. S18; Frazer, Belief in Immortality, i. 150-152.]
Was he cruel, warlike, passionate, generous, revengeful in the body, so will he be as a discarnate ghost. So, for instance, the New Guineans hold. Only account must be taken of a very common notion, that the ghost is endowed with increased power. One might find many reasons for this common idea. The general fearsomeness of the unknown and invisible, the fad that the ghost has joined the terrible host of free spirits, its very remoteness, combine to add the idea of power. That which is distant in space or time gains enchantment and enlargement from the imagination, which is the faculty most employed in this sphere. Australians credit to their ancestors deeds to themselves impossible, though they are themselves their ancestors reincarnate. The greed and liking for possessions which existed on earth are attributed in some parts to the spirit, and among the Bakongo, for instance, this desire is satisfied by placing all the deceased's wealth about the grave. The soul's assumed mobility,
[6. Neuhass, Deutsch Neu-Guinea, iii. 142 ff.
7. Roscoe, Baganda, pp. 282 ff.
8. Spencer and Gillen, Northern Tribes, pp. 489 ff
9. Weeks, Primitive Bakongo, p. 278.]
such as was displayed in its power to leave the body during life and to make investigations at even a considerable distance, is not lost but rather enhanced. It has become a free agent, no longer bound by the body's necessities and limitations of locomotion, at liberty to roam unfettered, to use in the wide universe its powers--those that remain or are acquired in its new condition. If it in earthly life could leave the body temporarily and like the lightning speed hither and thither, now, disfleshed, its mobility has gained by the change.
Especially is it believed that spirits acquire a larger knowledge. Not only do they gain a completer survey of the past and the present, but a knowledge of the future becomes theirs. According as their dispositions prompt, they become helpers of their survivors or hostilely active against them.
Particularly interesting in this connection is the relationship of the ghost and other beings to warning and prediction. Among the powers of the soul is that of return and manifestation to survivors. Melanesian, Andaman, and African ghosts, for instance, reappear to and converse with their people and become a medium of information. Particularly through dreams do they mediate--a performance recorded in antiquity and attested by present day belief over a large area. Indeed, it is through the dream that approach to human comprehension is most easily made by divine, superhuman, or discarnate powers, the spirit in this condition being loosed from fleshly trammels. The human spirit in sleep is regarded as not bound by quite the same inflexible laws to the bodily limitations. The employment of the dream as a means of information or warning at once occurs to the reader--Jacob, Joseph, Pharaoh, Nebucbadrezzar; clasctical cases will be found in Pindar, Olympiacs, XIII, 105 and Pausanias, X, xxxiii, II. It will be remembered that in an earlier section the importance of the dream as an index to animistic thought was dwelt upon at some length. One specimen of developed classical and philosophical thought on this has been summarized from Jamblichus.
[10. Seligmann, Melanesians, pp. 190 ff.; Klosts, In the Andamans, p. 296; Weeks, Congo Cannibals, pp. 264-265.
11. Herodotus, IV, 172; Pomponius Mela, I. viii. 8; Mauss, Origines des pouvoirs magiques, p. 15; Haddon, Anthropological Essays, p. 179.]
"There is nothing unworthy of belief in what you have been told concerning sleep and the meaning of dreams. I will explain it thus. The soul has a twofold life, a lower and a higher. In sleep the soul is released from the constraint of the body, and enters as one emancipated on its divine life of intelligence. Then as the noble faculty which beholds the objects that truly are, the objects in the world of intelligence, stirs within and awakens to its power, who can be surprised that the mind, which contains within itself the principles of all that happens, should in this, the state of liberation, discern the future in those antecedent principles which will make that future what it is to be? The nobler part of the soul is thus united by abstraction to higher natures, and becomes a participant in the wisdom and foreknowledge of the gods. Recorded examples of this are numerous and well authenticated; instances too occur every day. Numbers of sick by sleeping had their cure revealed to them in dreams. Would not Alexander's army have perished but for a dream in which Dionysius pointed out the means of safety? Was not the siege of Aphritis raised through a dream sent by Jupiter Ammon to Lysander? 'The night time of the body is the daytime of the soul.'"
The student of anthropology will at once recognize here the advanced justification for beliefs which go back very far for their origins. But even in the advanced stage of thought represented by Jamblichus there are present elements that are duplicable today in the most primitive regions.
Several doors open here to alluring bypaths--to inspiration, prediction, oracles, on the one side, these presuming a favoring: disposition on the part of the ghost; and, on the other, to necromancy and the "black art" or black magic, if the ghost or his control be evil. Melanesians and Africans say that the soul may return to seize and inspire the unconscious shaman or prophet to pregnant utterance. We have said "unconscious"--for it seems practically established that, in the earlier stages of culture, prediction and the delivery of the oracle took place only when the medium was in ecstasy. Vergil's description of the
[12. Theurgia or the Egyptian Mysteries, Part III, chap. vii.
13. Codrington, Melanesians, pp. 218 ff.; Roscue, Baganda, p. 113.]
raging sybil will recur to the classical student. Plato says that "inspired and true divination is not attained to by anyone in his full senses, but only when the power of thought is fettered by sleep or disease, or some paroxysm of frenzy." It is well known that the American Indians regarded the simple or mentally incompetent as peculiarly endowed and in closer touch with the supernatural than those possessed of all their mental powers. In the Old Testament there is an unconscious testimony to the veracity of many parts of the narrative, guaranteed by psychological conclusions, in the fact that the earlier phases of prophecy and prediction are described as involving the ecstatic state or a condition of unconsciousness. Such are the use of the dream, the case of Balaam, the prophets among whom Saul found himself, this form of affection being communicable or "catching"--compare the "dancing mania" of the middle ages-and Elisha, for whom music was in at least one case a prerequisite to the delivery of the oracle--the "hand of the Lord " (2 Kings 3:15) being the Old Testament expression for the modern psychological term
[14. Æneid, V1, 45 ff., 77ff.
15. Timæus, 71.]
"ecstasy" adopted from the Greek. So among perhaps most primitive peoples, like the Melanesians and Africans referred to above, warnings from the supernatural and even knowledge of other matters, as of charms, are supposed to be received under such conditions.
Ghosts do not figure merely as indicators of coming events or as guardians against evil fortune. Their larger capacity for action may make them powerful intercessors with still higher supernatural beings or spirits, through shamans who control them or know them intimately. Or their own success in their earthly vocation makes them interested in survivors who follow their trade. In Africa the spirit of a dead hunter is powerful to help in the chase, and is propitiated to that end. In Melanesia the help of ghosts in securing the right kind of weather, in performing feats of healing, in success with the fishing net or line, and in agriculture is obtained by sacrifices
[16. So the Australians: Howitt, Native Tribes, pp. 435-437. On the facts at large of Carpenter, Comparative Religion, pp. 181, 182.
17. Carruthers, Unknown Mongolia, i. 243.
18. Weeks, Primitive Bakongo, pp. 181-183.]
and offerings. Indeed, from the inhabitants of Ghosttown may come some of the good gifts, agricultural, for instance, which make life worth living. The spirits of the dead may keep a watchful eye upon survivors, preventing or punishing infractions of tribal customs that involve offence to themselves, and warning against repetition by inflicting sickness or failure in various enterprises. Foundation sacrifice had the purpose of procuring for the structure the protection of the spirits of the dead.
On the other hand, ghosts may be among the spirits whose malevolence needs to be guarded against. In fact, among the post mortem transformations may be that into ill disposed spirits. Usually, when this is conceived to be the case, the cause is found in some misfortune in life or death. Among the Ibo, for instance, a childless woman, a wifeless or moneyless man, or a suicide may as ghosts attempt to increase the population
[19. Codrington, Melanesians, pp. 132. ff.: Lambert, Murs et superstitions, pp. 24, 26, 218, 224 ff., 293 ff.; Turner, Samoa, pp. 345 ff.
20. Talbot, In the Shadow of the Bush, pp. 238-239.
21. Seligmann, Melanesians, pp. 192, 310.
22, B. D. Eerdmans, in Expositor, Nov. 1913, p. 197.]
of the underworld by attacks upon those left on earth. Similarly in New Guinea those who die in childbirth, suicides, and those who have lost their heads become maleficent. The Omahas hold that ghosts of the murdered return and inflict punishment by disease, or by causing the wind to blow from hunter to game and so to spoil his sport. Among Congo cannibals the soul seen in dreams is a wandering human spirit aiming at evil in its travels, and the witch doctor may be hired to kill it. The nostrils of the dead are plugged immediately after death to keep the spirit in the body as long as possible. If the ghost is for any reason unwelcome in the nether world and is driven out, it becomes malicious and aims at mischief, either inflicting positive ills by sending storms and like disasters or preventing success in various pursuits. In some cases ghosts are normally neutral, and their disposition and consequent actions depend upon the treatment they receive from
[23. Thomas, Anthropological Report, p. 312
24. Frazer, Belief in Immortality, i. 212.
25. Fletcher and La Flesche, 27th Report, p. 212.
26. Weeks, Congo Cannibals, p. 263-264, 269.]
the living. So that the well-being of survivors depends on propitiation by gifts and ceremonies or on manifestations of abiding affection. The duties of classic Greeks and Romans to their dead--careful and honorable burial, celebration by games at the funeral or on anniversaries--recur at once to the mind: and in these and other matters these peoples handed down in memory at least and often in ritual the doings and beliefs of far away ancestors. Close parallels to classic customs have been observed among African, Melanesian, and Polynesian peoples, where not only is the funeral offering placed on the ground, but dramatic performances in honor of the dead take place. Among some races, such as British New Guineans and the Mafulu, ghosts are always malevolent.
Among the exercises of the enlarged powers
[28. Williamson, South Sea Savage, pp. 65, 68, 74, 75, 76, 81 ff.; Roscoe, Baganda, pp. 116, 278, 286.
29. Taplin, Narrinyeri, p. 19; Curt, Australian Race, i. 87; Howitt, Native Tribes, pp. 461, 463, 473; Spencer and Gillen, Northern Tribes, p. 507, and Native Tribes, p. 511.
30. Talbot, In the Shadow of The Bush, p. 18; Brown, Melanesians and Polynesians, pp. 214 ff.; Milligan, Fetish Folk, pp. 233-236.
31. Williamson, South Sea Savage, p. 281 and Mafulu Mountain People, pp. 243 ff., 266 ff., 297 ff.; JAI, xxviii (1899), 216 ff.]
attributed to ghosts by quite diverse peoples is one which, as we shall see later, they possess in common with non-human spirits. This is the infliction of disease in an access of malignancy. Such a belief is held by American Indians, South Sea islanders, Hindus, New Guineans, and many others. They may inflict lockjaw by a blow, cause death, induce phthisis, and bring pestilence. Shamans and medicine men may use them to secure revenge or haunt the living; and this again calls up the need for exorcism. This gives rise to various devices and taboos, aiming at propitiating or deceiving the ghosts, such as change of names assigned to things belonging to the dead, or dropping out of the language words which contained the name borne in life, this going so far in some cases as to involve the destruction of huts, plantations, trees, and other possessions." It is quite in keeping with the
[32. Folk-lore, ii. 420 ff., 431; Kloss, In the Andamans, p. 305; Declé, Three years in Savage Africa, pp. 236, 344.
33. Talbot, In the Shadow of The Bush, p. 230; Weeks, Congo Cannibals, p. 266; Roscoe, Baganda, p. 100; Williamson, South Sea Savage, pp. 81 ff.; Crooke, Tribes and Caste, iii. 436.
34. Williamson, South Sea Savage, pp. 81 ff.; Roscoe, Baganda, p. 126.
35. Seligmann, Melanesians, pp. 631 ff.; Cambridge Anthropological Expedition, v. 250.]
whole conception of things that ghosts should be especially dangerous at night.
From all this, to anticipate slightly what is yet to come, fear of discarnate spirits may lead to a cult, a worship, which is apotropaic, deprecatory, or propitiatory in character. On the other hand, the sense of favors received or to come gives the rationale of a cultus which embodies more of gratitude and pleasure than of fear. With both these varieties of mental qualities attributed to ghosts, shared by them in common with non-human powers, it seems to require somewhat of ingenuity and a miscalculation or misappreciation of native human traits to force one to derive all worship from fear. Timor fecit deos is now hardly tenable in its original sense, in view of abundance of ascertained fads. Most of the animals, especially those domesticated, display amiable traits, including gratitude. We can hardly hold, therefore, that man, whether the product of evolution or of special creation, developed one of his noblest exercises, that of worship, from a sense of fear alone.
[36. Neuhass, Deutsch Neu-Guinea, iii. 64, 147.
37. Chalmers and Gill, Work and Adventure, pp. 84 ff.] |
Both the Crusades in the late eleventh century and the harsh persecution of Jews by the Bohemian king Vratislav II led to the migration of Jews from Bohemia, Germany and Austria to Hungary, where they found refuge; some of them settled in Slovakia. Under a "Jewish law" enacted by the Hungarian king Kalman, Jews were permitted to live only in cathedral towns and on bishops estates. In 1241 the Mongols (also known as Tartars) invaded Hungary, wreaking havoc and destruction. Jewish merchants made a major contribution to rebuilding the economy. A letter from King Bela IV refers to Jews in the cities of Pressburg (Bratislava), Senica, Trnava, Pezinok, Nitra, and Trencin. In 1251 the King granted the Jews a "privilege", a document promising them protection against attacks by Christians, a permanent legal status and other benefits. At the time, most Jews made their living in finance, a minority worked in commerce and the importing of goods, and others held positions in public administration or were involved in the minting of coins, to the displeasure of the Pope and Church leaders.
From the thirteenth century on the Jews were wards of the king and paid taxes to the royal treasury. They lived primarily in the cities, on separate Jewish streets allotted to them by the authorities in order to segregate them from the Christians. In some towns there was an organized community life and Jewish public institutions. Nitra, a major administrative and economic center in the Middle Ages, had a longstanding Jewish community that is considered one of the oldest in Slovakia. A document from the year 1113 mentions Mons Judaeorum in Nitra, where there was a Jewish cemetery. In the thirteenth century Jews also lived in the nearby suburb of Parovce, which was known as Castrum Judaeorum, i.e., "fortified Jewish settlement". Jewish sources also mention the Jewish community in Nitra; for example, in his book Or Zarua, written in the late thirteenth century, a leading halakhic authority, Rabbi Isaac ben Moses of Vienna, refers to a question asked by the community regarding marriage laws. At about the same time, Jewish refugees from Bohemia and Germany founded a community in Pressburg. In the fourteenth century the Pressburg community numbered close to 800 people and seems to have been the largest in the kingdom. In the late fourteenth century the spiritual leader of Trnava was Rabbi Isaac Tyrnau (the German name of Trnava), known as the leading Torah sage in Hungary and the author of Sefer Ha-minhagim, which describes the religious practices of the Jews of Hungary and the neighboring countries.
As the Christian population turned increasingly to religious extremism and social ferment in the fifteenth century, the Jews' situation grew worse. Anti-Jewish riots broke out in several places. In 1491 the authorities in Trnava spread a blood libel against the Jews, and twelve men and four women were burned at the stake on August 22 of that year. After the Turks defeated the Hungarians in 1526, the Jews were expelled from Pressburg, Trnava, and several other localities. In 1529 a blood libel was lodged against the Jews of Pezinok. Thirty members of the community were burned at the stake and the rest of the Jews fled the city. Mistreatment of Jews occurred in other communities as well. By the end of the sixteenth century the old communities of Slovakia had disintegrated, their members scattering in all directions. As a result of the severe persecution, the continuity of the community life in Slovakia was severely disturbed.
The immigration of Jews, beginning in the mid-seventeenth century and intensifying during the eighteenth century, gave rise to the Jewish communities of Slovakia that existed until the Holocaust. Some members of the Hungarian aristocracy realized the advantages that might be gained from the Jews economic activity and, unlike the townspeople, made efforts to encourage Jews to settle on their estates. The Jewish immigrants came mainly from the neighboring countries of Moravia, Poland and Austria. Jewish refugees settled in Nitra County in 1649 and later in the counties of Pressburg and Trencin. New edicts in Moravia and hardships in Poland spurred migration to Slovakia, and the area of Jewish settlement expanded northward and eastward. Most inhabitants of the regions in which new communities were established were Slovakian subsistence farmers, serfs of Hungarian feudal lords. The socioeconomic conditions created an interdependency between the two sectors of the population, especially in economic affairs. A census from 1746 shows that almost half the Jewish heads of household in Slovakia were natives of Moravia and Bohemia, 10 percent of Poland, 5 percent of Austria, and 35 percent of various locations in Slovakia or elsewhere in Hungary. By the middle of the eighteenth century there were several fairly large Jewish communities in Slovakia, some with populations in the hundreds. As Jews resettled Slovakia, an interesting encounter occurred between Jewish ethnicities and cultures. The Jews in the west were chiefly of the Ashkenazi type, who tended to be more educated and more open to influences of the surrounding culture and society. The Jews in eastern Slovakia, in contrast, followed Hasidic customs, spoke Yiddish, and resembled the Jews of Poland and Galicia in their way of life and dress. The two cultures were slow to blend, completing the process only in the second half of the nineteenth century.
The Jews' living conditions deteriorated drastically during the reign of Empress Maria Theresa (1740-l780). The regime promulgated various edicts and even threatened to expel the Jews from the empire. In 1749 the Jews were subjected to a special "tolerance tax" (taxa tolerantialis), a heavy burden for the immigrant families. However, their situation improved when Emperor Joseph II (l780-l790) instituted changes and innovations in the governance of the empire. In 1783 he issued an Edict of Tolerance (Sistematica Gentis Iudaicae Regulatio) for the Jews of Hungary, which referred mainly to place of residence, occupation, and educational matters. The Jews were granted permission to work in almost any occupation and to live in most parts of the empire, except in and around mining towns. Following the Edict of Tolerance, the areas of Jewish settlement expanded and new communities were founded. Jews became financially well off and in some places flourished spiritually and culturally. The Jewish community of Hungary also grew rapidly during this time, attaining a population of 83,000, about a half of them in Slovakia and the nearby Burgenland (Austria). Throughout most of the eighteenth century Jews lived in a relatively small area covering a few counties in the west and east. Large portions of central and northern Slovakia, defined as mining areas, were still off limits to Jewish settlement.
In the early nineteenth century the Jewish population of Slovakia grew quite vigorously. In the 1820s the region contained about one hundred organized communities, mostly in small towns and rural villages. The increase in the Jewish population and the gradual improvement in their economic condition were accompanied by a thriving religious life and the emergence of the first Torah centers. At that time there were seven relatively large yeshivas in Slovakia, headed by well-known rabbis and scholars. The most important of them was the Pressburg yeshiva, headed by Rabbi Moses Sofer (Schreiber), known as the Hatam Sofer, who was considered the leading halakhic authority of his day. In the midst of a liberalization wave in 1840, the Hungarian parliament passed several important legislative amendments pertaining to the Jews, mainly concerning places of residence and economic matters. Many Slovakian Jews moved to other parts of the empire around that time, and the old, traditional Jewish communities began to decline in number. New communities were founded in central and northern Slovakia, and the map of Jewish settlement in Slovakia attained geographic contiguity.
The "Spring of Nations" and the Hungarian uprising against Austrian rule (1848-1849) were not beneficial to the Jews. In March 1848 riots broke out in Pressburg (Bratislava) and nearby localities and spread to other regions. In many Jewish communities, especially in western Slovakia, houses were plundered and community institutions were destroyed; in some places there were casualties. Several localities were abandoned for some time.
In the second half of the nineteenth century, after the situation had stabilized, there were about 115 major Jewish communities in Slovakia with recognized rabbinical offices and 200 smaller communities with their own public institutions, chiefly in villages. Several dozen yeshivas, with thousands of students, became important Torah centers. Slovakian Jews played a significant economic role as mediators between the agricultural sector and the growing cities and were involved in trade in agricultural produce. They included merchants, tenant farmers and estate managers, but many worked in petty commerce as artisans or peddlers occupations that Jews in the Oberland had practiced for generations.
Whereas manifestations of socioeconomically motivated anti-Semitism that had erupted in the early nineteenth century declined when the revolution of 1848-1849 was suppressed, nationally based anti-Semitism increased over time. The Slovakian intelligentsia, which advocated a Slovakian national revival, regarded the Jews as tools of the hated Hungarian regime and, therefore, one of the causes of the plight of the Slovakian people. After Hungarian Jewry attained equal rights in 1867, anti-Semitic activity increased in Slovakia. Church representatives were actively involved in stirring up hatred of the Jews. In July 1882 some two hundred priests gathered in Topolcany to discuss the "Jewish question". The upshot of their deliberations was to urge the Hungarian parliament to revoke or limit the equal rights that it had granted to Jews in 1867.
From the beginning of the eighteenth century the empire had been undergoing a modernization process that was paralleled by similar changes in Hungarian Jewry. The ideas of the Enlightenment that seeped into Hungary during those years influenced many Hungarian Jews to cast aside their traditions and seek their future in European culture. Hungarian liberals encouraged Jews to abandon their traditional lifestyle and adopt the customs of enlightened Hungarian society. Attempts to amend or reform religious rules were met with vigorous opposition from both the religious leadership in Slovakia and a large portion of the Jewish population, which was characterized by a conservative way of life and a strong affinity for religion and tradition. The opposition to the reform trends was led, from the very beginning, by Rabbi Moshe Sofer of Pressburg. As a result of his vehement stance Pressburg became the center of the struggle against the Haskala (Jewish Enlightenment) movement and religious reforms. In the mid-nineteenth century the Kulturkampf within Hungarian Jewry intensified; and in 1865 dozens of Orthodox rabbis gathered in Michalovce, Slovakia. They adopted stringent resolutions including a boycott of synagogues that had instituted changes, and a strict injunction against sermons in German and Hungarian, secular education, and the study of foreign languages. The resolutions were given the force of a halakhic ruling, thereby further deepening the fissure among Hungarian Jewry and precluding all attempts to heal the rift. At the initiative of the authorities, representatives of Hungarian Jewry were invited to Budapest in December 1868 to promote the emancipation and to establish a countrywide umbrella organization for all Jewish communities in the country. The Liberal Jews welcomed the initiative, but the Orthodox regarded it with great suspicion. They feared that such an organization, supported by the authorities and the Liberal Jews, would diminish the communities' independence and undermine the status of the rabbi as supreme religious authority and sole arbiter. After all attempts at compromise failed, the Orthodox delegates walked out of the congress. The conflict between the camps reached its peak, causing a rupture between the communities and a deep polarization of Hungarian Jewry. As a result, two separate organizations of communities were established one for the Orthodox and one for the Liberals (Neologs). Some communities did not join either organization and retained their previous status; these communities were known as status quo ante. The Jews of Slovakia were more united in their affinity for religion than those in the rest of Hungary. Two-thirds of the communities, especially the old, traditional ones in the small towns and villages, joined the Orthodox organization. In parallel, most new communities, primarily in central Slovakia, joined the Neologs. In several locations two communities formed, one Orthodox and one Neolog, and they vied for control of community institutions and assets that had previously belonged to the joint administration.
On the eve of World War I the Jewish population of Slovakia numbered 140,000. Collectively they were quite diverse, comprising several religious factions whose members differed in their way of life, affinity for religion, origins, and cultural backgrounds. They spoke four different languages among themselves and with the Christians, belonged to various socioeconomic classes, and were divided among national groups. Most Slovakian Jews originated in villages and small towns and had traits characteristic of post-rural society.
After World War I social agitation grew in Slovakia. Economic hardship
increased following the return and demobilization of soldiers, and the
frustration was soon vented in several weeks of violent acts against Jews. The
burglary and looting of homes and businesses affected Jews of all classes, poor
and rich alike.
The new government adopted the Habsburg monarchies legislation regarding the Jews: namely, in the Czech province the Austrian laws remained in effect, whereas in Slovakia the Hungarian laws remained but with slight modifications. Attempts to establish a countrywide organization of communities representing all Jews in Czechoslovakia were foiled by the Orthodox leadership in Slovakia, which feared that this would undermine its hegemony among the religious Jews. Meanwhile, a new source of friction further exacerbated the tense relations between the factions. Under the Czechoslovakian constitution members of the Jewish religion could declare that they belonged to the Jewish nation. The Orthodox were vehemently opposed to defining Judaism as a national entity. To them Jewry was solely a religious community and would remain so until the Jews were redeemed from exile. Nevertheless, Jewish nationalism, represented by the Zionist movement, managed to establish itself and make inroads even among traditional Jews. Many Zionists, especially the young, came from Orthodox families; they joined Jewish national and Zionist organizations despite the opposition of the Orthodox religious leadership.
After Czechoslovakia was established the organizations of Jewish communities formed their institutions. The Organization of Autonomous Orthodox Congregations in Slovakia (OAOCS) included 170 of 228 congregations and about 75 percent of Slovakian Jews. Over time it became one of the most authoritative, powerful and influential Jewish organizations in Czechoslovakia. Meanwhile, Agudath Israel, which had close ties to the OAOCS, complemented it by carrying out a wide range of activities in dozens of branches. Agudath Israel focused primarily on individual life and on education consistent with Orthodoxy. The ties between the two organizations, strong to begin with, became even stronger when Rabbi Samuel David Ungar, one of the leaders of the OAOCS, was named president of Agudath Israel in Czechoslovakia. Agudath Israel and its youth movements engaged in ramified social and welfare activities and educational initiatives, including preschools, a Beth Yaakov school system for girls, and camps for children and teenagers. At the end of World War I there were twenty-nine Neolog communities in Slovakia. Cut off from their center in Hungary, they encountered complex organizational problems that threatened their survival. In 1925 the Neologs decided to form a joint organization with the fifty-five Status Quo Ante communities, which since 1928 had been known as Jeshurun.
In 1920 there were seventy-seven Jewish primary schools in Slovakia and two high schools. Forty-six of the schools were Orthodox; the rest were liberal. Noteworthy, only 45 percent of Jewish children attended Jewish schools; the majority attended public schools. Due to the growing percentage of students in schools that taught in Slovakian, the language spoken by the young Jews changed. Instead of German and Hungarian, the main languages used by Slovakian Jews in the past, the young people adopted Slovakian as their vernacular. Alongside the formal educational system, an extensive system of Torah education institutions functioned in Slovakia, including boys' schools (cheder and Talmud Torah), Beth Yaacov seminaries for religious girls, and more than thirty yeshivas that were accredited as religious educational institutions. In 1930 the Pressburg (Bratislava) yeshiva and three other yeshivas were accredited by the Czechoslovakian Ministry of Education as institutions of higher education.
Vibrant Zionist activity was enjoyed in inter-war Slovakia and new branches of the movement were established in dozens of cities and towns. In addition to the Zionist Organization, there was a Jewish National Party (Zidovska Strana) that focused on domestic Jewish policy and representation of the Jews national, religious, economic, and social interests vis-a-vis the authorities. The party did well in elections for town councils and county and regional assemblies. Although the Orthodox establishment doggedly fought Zionism and the Jewish National Party, the Zionists penetrated even the Orthodox communities. In the mid-1930s about 20 percent of Slovakian Jews purchased the Zionist "Shekel", thereby becoming formal members of the Zionist Movement.
Zionist youth movements predated World War I. In 1919-1920, after the wartime turmoil in Slovakia had waned, Zionist youth groups reorganized and opened clubs in all the major cities and some small towns. In 1924, representatives of Zionist youth groups met in Nitra and together founded He-haluts. Religious teenagers established Mizrachi Youth. Organizational and social activity in the Zionist youth movements reached its peak in the 1930s, with tens of thousands of young members from all strata of Jewish society. Hundreds of them emigrated to Palestine after undergoing training for this purpose in Slovakia. The Maccabi sports association launched social and cultural activity with Zionist national leanings. Thousands of youngsters and adults, Orthodox and Liberal alike, were active in clubs throughout Slovakia.
During the period of the Republic of Czechoslovakia, Slovakian Jewry did not undergo any major demographic changes. In the 1930s about 140,000 Jews lived in 2,337 localities. Some 55 percent of them lived in villages and small towns with populations of less than 5,000; 18 percent lived in the two largest cities (Bratislava and Kosice), and only 25 percent lived in cities of l0,000-25,000 residents. The socioeconomic status of Slovakian Jews also remained constant during this period: 72 percent of Jewish breadwinners were self-employed, 15 percent were wage earners, and 13 percent were practitioners of liberal professions, white-collar workers, and brokers.
In the second half of the 1930s, as political tensions mounted in
Czechoslovakia, anti-Semitic sentiments increased among large segments of the
Slovakian people. The nationalist parties that agitated for Slovakian national
autonomy, especially the Slovak Peoples Party, held strong anti-Semitic
positions and incited against the Jews. On October 6, 1938, after the Munich
agreement in late September had forced Czechoslovakia to cede territory to the
Third Reich, the Slovak Peoples Party declared extensive autonomy in Slovakia
and instituted a one-party totalitarian regime. On November 2 large portions of
southern Slovakia home to more than 45,000 Jews were annexed to
Hungary. The Slovakian government blamed the Jews regarding their supposed
support of the annexation and started removing thousands of Jewish families who
held foreign citizenship to the Hungarian and Polish borders. Few of the
deportees were permitted to return to their homes.
While the deportations were in progress, an underground cell known as the Working Group, headed by Rabbi Michael Dov Weissmandel and Gisi Fleischmann, formed within the Center of Jews in an attempt to stanch the deportations and aid the deportees. The Working Groups success in halting the deportations from Slovakia encouraged its members to intensify their efforts to save all of European Jewry by negotiating with the Nazis. The resulting initiative, devised by Rabbi Weissmandel, was known as the Europa Plan. In April 1944 two Slovakian Jews, Alfred Wetzler and Walter Rosenberg (Rudolf Vrba), escaped from Auschwitz. The Working Group took detailed testimony from them about the nature of the Auschwitz camp and the extermination methods used there. A sketch of the extermination facilities, based on the description given by the two men, was attached to their testimony. The Auschwitz Protocols, as the document was known, were sent from Slovakia to the free world, where they reached Jewish organizations. The Working Groups purpose in releasing the Protocols was to sound an alarm and marshal worldwide public opinion in favor of extensive rescue activities and bombing of the railroads leading to Auschwitz and the extermination facilities themselves.
Jews joined the first resistance groups that organized in Slovakia in 1942. Most of them were active in the Communist and Czechoslovakian underground movements. A Jewish underground group was also formed in the Novaky labor camp in 1942. These groups tried to prevent a resumption of the deportations and started preparing for active resistance. On August 29, 1944, an armed uprising broke out in Slovakia in an attempt to overthrow the pro-Nazi regime and reestablish the Republic of Czechoslovakia. As the partisan attacks intensified, the German army invaded Slovakia. The Jewish group from the Novaky camp was assigned to halt the advance of the SS troops along one of their main attack routes. The German invasion of Slovakia augured ill for the last remaining Jews. Many tried to reach the rebel-controlled area in the hope of surviving, and thousands of destitute Jewish refugees gathered in the rebels main stronghold, Banska Bystrica. Parachutists sent by the Yishuv (the pre-Israel Jewish community in Palestine), headed by Haviva Reik, mobilized to help them. Although the uprising was quelled two months after it began, rebel units, including Jews, fled to the mountains, where they continued to engage the enemy until the liberation. About 1,600 Jews fought in various partisan units; approximately 170 of them were killed.
After the uprising was suppressed, the Germans took over the authority for Jewish affairs. SS officer Alois Brunner, one of Eichmanns assistants, went to Slovakia to deport all Jews irrespective of their status or their "certificates of exemption". Those Jews who were captured by the Nazis and their Slovakian accomplices were taken to the Sered camp. The deportation of the remaining Jews in Slovakia resumed on September 30, 1944. From then until March 31, 1945, some 12,000 Jews were deported from Slovakia; only half survived. Another 2,500 Jews were murdered on Slovakian soil during this period. Additional victims among Slovakian Jews were those who had fled to Hungary and were deported from there to the extermination camps.
The Jews in the territories annexed to Hungary in 1938-1939 met the same fate
as those in the rest of Hungary. After the annexation of these territories,
Hungarians began to persecute the Jews and accused them of supporting
Czechoslovakia. Of the 10,600 business owners only 4,500 were permitted to keep
their establishments going. Beginning in 1940, close to 7,500 men from southern
Slovakia were taken to work in labor battalions; few survived. Several thousand
Jews lacking Hungarian citizenship were deported in 1941 to the occupied part
of Ukraine, where most of them were murdered. After the Germans occupied
Hungary (March 19, 1944), new anti-Jewish edicts were promulgated.
Ghettoization of the Jews began in the second half of April 1944. The first
transports from the territories annexed to Hungary left for Auschwitz in the
second half of May 1944; the rest of the deportations occurred in June of that
year. Of some 45,000 Jews who lived in those territories, 10,000 survived.
About 100,000 Slovakian Jews 73 percent of their number in 1938
perished during World War II.
JewishGen, Inc. makes no representations regarding the accuracy of
the translation. The reader may wish to refer to the original material
JewishGen is not responsible for inaccuracies or omissions in the original work and cannot rewrite or edit the text to correct inaccuracies and/or omissions.
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- Consul (abbrev. cos.; Latin plural consules) was the highest elected office of the Roman Republic and an appointive office under the Empire. The title was also used in other city states and also revived in modern states, notably in the First French Republic. — “Consul - Wikipedia, the free encyclopedia”,
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- Représentants des collectivités locales De g à d B Casaurang CG 30 J P Boré Région D J Valade Ville Consuls généraux De g à d Espagne Italie Allemagne
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- Saint Antonin Noble Val Tarn et Garonne Maison des Consuls which housed a court of justice ferocious mutual beard pullers on a capital The best of all from the ruined Abbey of La Sauve Majeure Gironde and now in the Cloisters
- Les Consuls
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- The Maison des Consuls Mirepoix proposes you its 8 rooms Secured online booking La Maison des Consuls is situated in a beautiful and widely visited medieval square in the heart of Mirepoix
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- with no separate chassis and MacPherson front struts were used for the first time on a British car It had a live axle with leaf springs at the rear hydraulic drum brakes front rear The Consul had a 1508cc 4 cylinder engine which developed 48bhp mated to a 3 speed column change gearbox with synchromesh only on second and top gear The handbrake was operated by a pull
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- Nous allons voir la cathédrale Saint Maurice qui possède la nef la plus large de France Sa flèche de pierre est très belle mais nous apprécions aussi les très nombreuses gargouilles
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- Une délégation conduite par S E l ambassadeur Richard Zady a pris part aux travaux de la première Conférence des Consuls Honoraires de la Cote d Ivoire du 17 au 22 septembre Palais des
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- Klik voor meer foto s Consuls zijn vrijwilligers die inzetbaar zijn voor allerlei verschillende taken op de werkterreinen van de ANWB Zij spelen al sinds de oprichting van de ANWB een
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- Consuls From the Capitoline Hill Museum
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- XIIème siècle le péché et le châtiment du péché La rédemption n est pas représentée la justice n étant pas concernée par ce concept Sur la gauche de l image on peut voir Adam et Ève vus en gros plan sur la page précédente symbolisant l acte de pécher
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- Mr Howard James Pym Honorary Consul Of Latvia attended the 4th meeting of Honorary Consuls in Riga on the 3rd and 4th of July Latvias Consuls of the World Mr Howard J Pym Honorary Consul of Latvia Pictured above 1st from right 2nd row together with His Excellency Ivars Godmanis Prime Minister of Latvia pictured
- la royauté Le Conseil d État eut une activité très importante sous le Consulat et le Ier Empire 1799 1814 C est à lui que l on dut notamment la préparation des codes napoléoniens
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- moustachus étant punis par les aigles symboles du pouvoir divin d autres pécheresses leurs tresses postiches faits avec les cheveux de religieuses et d indigentes saisies par les becs des aigles aux
- Other capitals show a sinful woman in the clutches of a monster though in the context of the Maison des Consuls it might be Justice triumphing over Injustice a pair of mutual beard pullers symbol of strife
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- University of Pangasinan: Consuls of Goodwill Faculties/CI's Dance Presentation bonggah! lots of yells! watch it! wink!
- British Ambassador & consuls chat to Phuket press Comments by new British Ambassador, Asif Ahmad, on first visit to Phuket; recorded at Millenium Hotel, Patong, on Sunday 27 February 2011 Video starts with his comments on free trade area in ASEAN; then new requirement for UK visa residency applicants who are spouses of British citizens having to pass an English language test before entry is given; then on safety issues for tourists, jet skis and tuk tuks, to be raised with the Phuket Governor in a meeting today; the quarterly meeting of Hon Consuls with Governor; declining comment on Thai PM's possible dual nationality status; Thai visa rules on foreign resident's applications; help for Thai police who lack English skills; (after about 15 mins) British Hon Consul Martin Carpenter joins in with comments on tourist & normal police assistance; his contact details not currently being on Embassy website; any upgrading to consulate level & requests for visa applications in Phuket; police help in accidents; (after 20 mins) Ambassador said Consular staff would soon visit Phuket for temporary services 'surgery' (after 21 mins) British Consul Micheal Han*** comments on UK staff being cut back and replaced by local staff; decision on whether to upgrade to a consulate here; (after 23 mins) Hon Consul on the number of British residents, prisoners, and tourists in Phuket & region; assistance in other southern provinces (after 27 mins) Ambassador on Locate system for British residents to register on but increasing use of social media ...
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- MFA met Hon Consuls in Phuket pt 1 of 2 (raw video footage of beginning of meeting) International diplomatic representatives in the Andaman region met with senior officials from Ministry of Foreign Affairs and Phuket authorities to lay out clear channels of communication, strengthen relations and enhance mission efficiency It was a busy day on Monday for Consul-Generals and Honorary Consuls representing at least 18 foreign countries as they attended two important meetings in Phuket, both gearing towards improving communication channels among themselves and the Thai authorities. In the morning about 17 Consul-Generals and Honorary Consuls, who are mostly covering Phuket, Phang-Nga and Krabi, had a friendly session with officials from the Ministry of Foreign Affairs' Protocol Department at the Adamas Resort and Spa on Nai Yang beach. The session was chaired by the Chief of Protocol Department, Mr. Bansarn Bunnag. German Ambassador to Thailand, HE Dr. Hanns Heinrich Schumacher was also present. The meeting discussed clear channels of communication, both in normal times, but with contingency plans in case of emergencies, following international practice. The meeting was shown what's called a 'Phuket Model' design with concise channel of contacts and communication in the wake of any needs to handle stressful situations. The Thai Chief of Protocol told the international representatives that the ministry designated the Chief of the Passport Office in Phuket, Thanawat Sirikul, as the first contact person in the Andaman ...
- CONSULS AND YOU.mov Elihu Burrit Library instructional video about Consuls.
- Amadou Ballaké et Les 5 Consuls - Ligda remba. 45 rpm disc : Amadou Ballaké et Les 5 Consuls - A la mémoire du regretté Demba / Ligda remba. Amadou Balaké was born in 1944 in Ouahigouya, in the northern-east of Burkina Faso. After the death of his father, he relocated with his mother in the capital of the country, Ouagadougou, where he began to meet several local musicians. Shortly after he went to Mali to work as a chauffeur apprentice, and only returned to Ouagadougou six years later. He worked some time in construction, and then became a taximan. After the loss of his vehicle, he went to Mali and was engaged as a professional musician in the Grand Hotel orchestra. In 1963 he left Mali to play in the Tropicana orchestra in Abidjan, capital of the Ivory Coast, but only six months later he went to Guinea. Leader of the Bafinf Jazz band, he participated there in many musical competitions, and regularly played for president Sékou Touré meetings across the country. He returned to Burkina at the end of the 60's, and became very popular in the country thanks to his work with the Harmonie Voltaïque band. Convided by the Nigerian producer Aboudou Lassissi, he went back to Abidjan at the beginning of the 70's .
- Balitang America: Opera Consuls Don Tagala meets up with French Consul Members who are trained by Filipinas to take the stage.
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- MFA met Hon Consuls in Phuket pt 2 of 2 (raw video footage of beginning of meeting) International diplomatic representatives in the Andaman region met with senior officials from Ministry of Foreign Affairs and Phuket authorities to lay out clear channels of communication, strengthen relations and enhance mission efficiency It was a busy day on Monday for Consul-Generals and Honorary Consuls representing at least 18 foreign countries as they attended two important meetings in Phuket, both gearing towards improving communication channels among themselves and the Thai authorities. In the morning about 17 Consul-Generals and Honorary Consuls, who are mostly covering Phuket, Phang-Nga and Krabi, had a friendly session with officials from the Ministry of Foreign Affairs' Protocol Department at the Adamas Resort and Spa on Nai Yang beach. The session was chaired by the Chief of Protocol Department, Mr. Bansarn Bunnag. German Ambassador to Thailand, HE Dr. Hanns Heinrich Schumacher was also present. The meeting discussed clear channels of communication, both in normal times, but with contingency plans in case of emergencies, following international practice. The meeting was shown what's called a 'Phuket Model' design with concise channel of contacts and communication in the wake of any needs to handle stressful situations. The Thai Chief of Protocol told the international representatives that the ministry designated the Chief of the Passport Office in Phuket, Thanawat Sirikul, as the first contact person in the Andaman ...
- Consuls General iftar Dinner,August 2010
- Little Caesar & the Consuls - (My Girl) Sloopy One of many versions of this song (originally by the Vibrations), this was a #1 hit in Canada in July of 1965. It peaked at #50 in the US. This version is slower and more emotional than the McCoys more pop oriented version (retitled Hang on Sloopy) which went to #1 in the US later that same year.
- (Part 1 of 2)Call by the Delegates of the 5th Ambassadors/Consuls General and Tourism Directors Rizal Hall, Malacanang Palace July 11, 2009
- Jump Jeremiah - Mike Ford & the Consuls British group formed in the 1950s - my handsome and talented nephew, Geoffrey, was part of the group.
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- The evolution of gaming consuls Most of the popular consuls and their prices at retail. Final multimedia project. You'll notice i left out the Wii and computer systems. sorry i was time limited when i made this.
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- Consuls Meeting Following the success of the honorary consuls' meetings in Phuket, the Thai Ministry of Foreign Affairs has announced that a nationwide gathering of al honorary consuls will be held at the ministry in Bangkok next month.
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- 24 February 2010 The Society of Foreign Consuls in NY NYSE Euronext Opening Bell. The Society of Foreign Consuls in New York rang The Opening Bell.
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- CLANS YOUTUBE SUBSCRIBE FOR 10TH LOBBY AL CONSULS SUBSCRIBE
- 5th Ambassadors Consuls General and Tourism Directors Tour of the Philippines Join us as we embark on another unforgettable trip to the Philippines: The 5th Ambassadors, Consuls General, and Tourism Directors Tour (ACGTDT) of the Philippines. The one-of-a-kind tour takes place on July 9-17, 2009. For only $1399, inclusive of round trip international airfare, to and from Manila and San Francisco, Los Angeles, or Las Vegas, 3-night hotel accommodations at the Dusit Hotel in Makati with daily breakfast, Dinners with cultural shows and entertainment, city tour of Manila, choice of day tour to famous out of town destinations, business opportunities activities, a visit and lunch at Malacanang Palace with tour of the museum, and a photo opportunity with Her Excellency Gloria Macapagal-Arroyo. Log on to www.experiencphilippines.ph for more details.
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- Penang set up Council of Foreign Consuls Komtar, June 1st 2010 - Penang has set up the Council of Foreign Consuls which consist of 19 foreign Hon. Consul, Consul-General and other representatives. At the same time, CM Lim Guan Eng also introduced Khoo Kay Peng and Liew Chin Tong as two of several members who will represent Penang in Kuala Lumpur to promote Penang with all the diplomatic offices.
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- Sri Lankan Consuls urged to counter ill informed charges 19/01/2009 President Rajapakse has called on Sri Lankas Consuls abroad to inform the world of the achievements made towards restoring democracy and freedom in the North. They should do so and counter the ill informed charges being made about the plight of the innocent Tamil people there..
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- Governor met Hon Consuls in Phuket (raw video footage of beginning of meeting) International diplomatic representatives in the Andaman region met with Phuket authorities to lay out clear channels of communication, strengthen relations and enhance mission efficiency (story continued from 2 videos of morning meeting) In the afternoon, the group had a meeting with the Phuket authorities, chaired by Phuket Governor Wichai Praisa-Ngob at the Provincial Hall. The Chief of Protocol Department, MFA, his deputy and German Ambassador to Thailand were again present. The governor told the foreign envoys that he will make sure that his Thai officials are efficient and friendly to welcome tourists and visitors starting from the airport, emphasizing safety as a main priority for visitors and residents. Apart from establishing concise channels of communication to facilitate work efficiency, the focal issue in the afternoon centred around cases concerning foreigners in Phuket and how people involved are treated. The Provincial Police Commander Pikad Thantipong reported the number of cases concerned with foreign nationals. It was agreed earlier that the ambassadors, consuls or representatives be notified in writing when such cases occurred. Pol. Major General Pikad however admitted that there may be delays in some cases. He pointed out the comparative figure of fatal cases for foreigners in Phuket, stating that in 2008 there was 61 cases, in 2009 there was 28 cases and this year up to August there was 52 fatal cases. In ...
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“Archived from groups: rec.games.miniatures.warhammer (More info?) Howdy, with the release of 4th, I've divorced the Angels of Vengeance I had from the DA codex, and turned them into Black Consul”
— Black Consuls, force review request - Games-General,
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— : Blog Articles " Mexican Consuls Increase Budget,
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— History Blog " Roman consuls,
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“Re: Update on your consuls. hadji hadji@- New provinces II "Antonio Grilo" amg@- Re: The Plebeian cults of Ceres, Liber and Libera. jmath669642reng@-) Alexander I.C.P.M.. Claudia Aprica Re: Update on your consuls. SFP55@- Re: Praefecti/Hispania”
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“Groups: Accreditation Editors, Railway Construction, Regional Consuls, The GrapeVine, Turner & Townsend Groups: Brazil, IT Engineering Projects, Regional Consuls. Hello, I think that besides the forum, we need an agenda. We should”
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— Hotel Des Consuls | Gordon Frickers' Blog,
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Who's most likely to get Parkinson's disease?
In the U.S., the disease is more likely to affect men than women, more likely to affect white or Hispanic individuals than black or Asian people and much more likely to affect those over sixty than those younger, according to a study published in the American Journal of Epidemiology in 2003.
Other research suggests that globally, you're more likely to get Parkinson's if you live in an industrialized country like the U.S. than if you live in a less developed nation. However, there's no clear-cut data on the demographics.
"It's hard to believe that we don't have that data now," says Dr. J. William Langston, founder and director of the Parkinson's Institute, "but ... it's not available anywhere in the world." In part, this is because Parkinson's disease affects a relatively small population -- at most one to two percent in the United States -- and because of difficulties in consistently diagnosing the disorder.
Is the risk of getting this disease growing?
There's even less data on how the rate of Parkinson's cases have changed since it was first identified in 1817. The number of cases is rising as the global population ages, but it's not clear whether Parkinson's is affecting a larger percentage of the population than it has in the past. Some research has found that the risk of getting Parkinson's in the U.S. rose slightly between 1935 and 1985, but when it comes to tracing the history of the disease, major data holes remain.
"Certainly, there have been individuals described hundreds and hundreds of years ago that have what sounds very much like Parkinson's disease," says Dr. Michael Zigmond, a Parkinson's researcher at the University of Pittsburgh. However, he adds, "We can't easily evaluate people who are not here now, even through the eyes of a neurologist 50 or 100 years ago." What we call Parkinson's today, doctors might have diagnosed differently in the past. What doctors used to call Parkinson's we may now have identified as something else.
However, there are some new projects underway on the demographics of the disease. For example, the California Parkinson's Disease Registry will collect detailed information on every Parkinson's diagnosis in that state. Says Dr. Langston, "This will tell us for the first time, is the disease increasing, changing with time? Are there pockets or clusters of the disease? Are there differences in rural versus urban areas, socioeconomic differences, etc.?"
What causes Parkinson's?
Even though Parkinson's was formally identified almost 200 years ago, scientists are still trying to find out what causes it. "We still don't have a smoking gun, that's for sure," says Dr. Langston. "But that's what we're looking for."
The greatest obstacle to pinpointing what causes it is that there seems to be no single cause. In some cases, the disease appears to be a genetic defect, while in others, exposure to toxic substances or certain viruses seems to be a factor. Some scientists believe a mix of environmental exposures and underlying genetic sensitivities may ultimately explain what triggers the disease.
What about environmental factors?
Sorting out the possible environmental influences on Parkinson's disease is a big challenge.
In the early 1980s, after discovering the toxin MPTP -- very similar to a compound found in a common pesticide -- could induce Parkinson's virtually overnight, scientists expected to track down the chemicals responsible for causing the disease in a matter of a few years. Nearly three decades later, "Very few individual specific compounds have been identified," says Dr. Langston. "Most of these are modest risk factors, two to threefold increase in risk."
Scientists have tracked disparate leads in search of an environmental explanation. They've researched Parkinson's-like symptoms in an individual who ingested petroleum products, in Taiwanese women who'd contracted a herpes virus and in people from Guam who eat a type of seed known to contain a neurotoxin. While finding a common thread among these scenarios has proved nearly impossible, much of the research has led back to pesticide exposure as the leading suspect in causing a large number of Parkinson's cases.
How clear is the link between pesticide exposure and Parkinson's?
According to Dr. William Langston, "I don't think we're yet at the point of being able to say unequivocally that if you live in an area where there's more pesticides than other areas, you're at a higher risk." But new research may illuminate the link.
One study is mapping Parkinson's cases in California against areas known to have high pesticide use. Another is collecting data from over 50,000 licensed pesticide applicators in two states to track whether these individuals face an increased risk of getting Parkinson's. "This, I think, may really answer the question of not only whether pesticides are an increased risk, but also, specifically, which pesticides," Dr. Langston hopes. "Because it's their profession, they know what they use, as opposed to, I couldn't tell you what I sprayed with in the garage last week."
Is there a gene, or genes, that causes Parkinson's?
So far scientists have identified six forms of genetic Parkinson's and are searching for more. However, the number of cases with a genetic link is a tiny portion of the overall number of cases. Still, doctors hope that better understanding the genetic form of the disease could help unravel the mysteries about other forms of Parkinson's.
But the genetic component of Parkinson's has also proved more complicated than once thought. Scientists have determined there's no single gene for Parkinson's -- one gene may trigger the disease in one family, while a different gene triggers it in another family.
Scientists also suspect that in some cases, developing genetic Parkinson's may require having multiple trigger genes, because studies of families with a high incidence of the disorder have shown that some individuals can carry the main gene identified as causing Parkinson's yet never develop symptoms.
How far along are scientists in figuring out the interaction -- the connection -- between genes and the environment?
A large study of over 20,000 identical twins -- who share the same genetic code -- is underway to trace when different environmental exposures can trigger Parkinson's in one twin but not the other and under what circumstances both twins develop Parkinson's.
Why is Parkinson's so hard to diagnose?
Diagnosing Parkinson's can be a challenge in part because some mild Parkinson's symptoms at first just seem like the universal effects of aging -- a tremor in the hand, difficulty balancing and shuffling the feet.
"I think everybody gets a little Parkinsonian as you get older," says Dr. Clive Svendsen, who is studying stem cell therapies for Parkinson's at the University of Wisconsin. "And in fact, most of the literature points out a gradual reduction in dopamine neurons" -- the neurons whose death causes Parkinson's best-known physical symptoms -- "over time in everybody."
What are the hallmark symptoms? Do all patients have the same symptoms?
There's no clear epidemiological record of whether the specific symptoms doctors look for and patients express have evolved over time, but the hallmark of Parkinson's has always been tremors throughout the body. Making a definitive diagnosis has often proved difficult, however, because the exact nature of the symptoms can vary widely from patient to patient. Recent studies show that even members of the same family who share a single genetic form of Parkinson's may display very different symptoms. One patient may have foot tremors accompanied by difficulty sleeping, while another may have hand tremors and difficulty keeping his balance.
What are the new discoveries about Parkinson's symptoms?
In recent years scientists have found that Parkinson's is much more than a disease of shaking limbs.
"When I started my residency, this was a very simple disease," recalls Dr. William Langston. "A number of cells die in a small area of the brain that made a chemical called dopamine. When they died, you had no more dopamine. Without dopamine, it's difficult to move. ... And that's the way we diagnosed it. When dopamine's down, you got rigid, you developed a tremor, gait became slowed and shuffling, etc. Any neurologist can diagnose that."
But now, says Langston, "At this point in time, we know that Parkinson's is a much more complicated disorder. Many different areas of the brain can be affected. It probably evolves in a very specific order, starting in the low brain stem and then eventually affecting other areas, including the nigra, which causes Parkinsonism [the tremors]. But all of these other areas of the brain that are affected can also cause symptoms."
These newly-recognized symptoms range from loss of the sense of smell to digestive problems to depression.
Why has it taken doctors so long to recognize the wide array of symptoms now connected to Parkinson's? "Because [patients] don't come to neurologists if they have, say, sleep disorders or loss of sense of smell or even constipation, which is a very bothersome symptom in Parkinson's," says Langston.
How close are we to having an earlier, more accurate diagnosis of Parkinson's?
The National Institute of Neurological Disorders and Stroke, part of the National Institutes of Health, is working on improving its guidelines for diagnosing Parkinson's to reflect the latest science. Other researchers are focusing on finding new ways of identifying the disease altogether. Some are searching for a Parkinson's biomarker -- a biological trait displayed only in people who have or are at risk of developing the disease.
And now that sense of smell has been tied to Parkinson's, a study led by the Institute for Neurodegenerative Disorders and the University of Pennsylvania is seeking to determine if how well a person performs on a test to identify 40 common odors can predict the individual's likelihood of developing Parkinson's. So far, they've found the average person can identify 35 of the odors, while the average diagnosed Parkinson's patient can only correctly name 20.
How close are scientists to finding a cure?
As scientists learn more about the great complexity of Parkinson's disease, hopes for finding a cure within the next few years is fading. But all the new information is paving the way for inventing better treatments that won't cure Parkinson's completely but will minimize the disease's effects.
"What we have to think about is whether the patient would be happy with still having constipation, still having other side effects of the disease, but being able to maintain a movement without arresting tremor and being able to initiate their own movement," says Dr. Clive Svendsen. "We're looking for a treatment that ... isn't curing it, but it's making the quality of life better."
If there's no cure in sight in the near future, what are the treatment options?
Since the 1960s, the leading treatment for Parkinson's has been the drug Levodopa, or L-dopa, a compound that can counteract the loss of dopamine neurons that causes the tremor symptoms of the disease. In the past couple of decades, deep brain stimulation -- a treatment that involves implanting a pacemaker-like device to deliver electric shocks to the parts of the brain that are damaged in Parkinson's patients -- has become another widely-employed treatment that generally produces positive results.
But the effectiveness of L-dopa decreases the longer a patient uses it, and deep brain stimulation requires a highly invasive operation, so scientists continue to look for better Parkinson's treatments.
As explored in My Father, My Brother and Me, doctors have recently started studying the power of exercise as a therapy to not only keep Parkinson's patients physically healthy but beneficially alter their brain chemistry.
What are some of the newer treatments being pursued?
Another area that offers great promise is neuroprotection, explained by Dr. Langston as "the idea of protecting nerve cells from dying or damage, ... another great Holy Grail in the field of neurology."
Scientists studying neuroprotection are looking for chemical substances that, when introduced into the brain, could either blunt the effects of Parkinson's disease on a patient after he is diagnosed or, better yet, prevent Parkinson's altogether. "The problem is," according to Dr. Langston, "I think it's a laboratory concept that jumped into the clinical all too quickly. In a laboratory animal, you can actually measure nerve cells after an experiment. We can't do that in humans. We have no way to really show we've slowed down the progress of cell death in living humans."
With thousands of potentially helpful chemicals to test, and without the scientific ability to widely screen substances for their ability to protect human brains against damage, Langston suggests researchers focus less on the effects of these substances on the brain itself and more on their effects on the patient as a whole.
"For the moment, where I think we need to focus in clinical trials is delaying disability. We clinicians can measure that. So start a patient on your neuroprotective agent. We can't prove [the neuroprotection], but if disability is really delayed or completely stopped, I think that would be very compelling."
Other scientists are working on refining the mode of delivery for a known neuroprotective substance. "There are some very powerful drugs for Parkinson's disease that are very difficult to get into the brain," says Dr. Svendsen. "GDNF [glial cell derived neurotrophic factor] is one of those drugs. ... Even at late stage Parkinson's, you have lots of neurons left, they just don't have any dopamine left in them, and they're very shrunken. And in that sense, [growth factors like GDNF] might again be like the fertilizer. Put it onto those cells, and even though they've practically disappeared, the growth factor will make them rejuvenate and start to produce dopamine again."
The problem, says Svendsen, is, "You can give [GDNF] peripherally in the blood, but it doesn't penetrate the blood-brain barrier [to get into the brain]. The idea is to design stem cells that make this drug, put the stem cells in the brain, and then they'll deliver it -- rather like a Trojan Horse. The brain accepts the stem cell because it's [a brain cell] and it's going to integrate and migrate and get into the brain tissue. ... And that's how we can sneak drugs into the brain, [overcoming] this powerful blood-brain barrier, which usually blocks this process."
The major obstacle Svendsen needs to solve to implement this type of therapy is creating an off switch of sorts for the cells that will produce GDNF. "It does remind me of the Walt Disney Fantasia movie and 'The Sorcerer's Apprentice,'" he says, "when Mickey Mouse is in the basement with the sorcerer and he gets a great way to make the broomsticks carry the water up and down. And, of course, it goes horribly wrong because [they] carry too much water, and it's overflowing, [so he] chops the broomsticks in half [but] they just keep carrying the water up and down.
"I wake up at night with that dream in my head, thinking, boy, if we put cells in the brain that produce [GDNF] that we think it's great and then it has a toxic effect, we can't switch it off. ... But on the other hand, we need to move forward in these diseases. ... It's a difficult one."
Could Parkinson's be a "gateway disease?" Could solving it clear the way to understanding -- even curing -- other neurological conditions?
"I think there's a general sense in the scientific and medical community that solving any of these major diseases -- Parkinson's, Lou Gehrig Disease, Alzheimer's -- could have an enormous impact on the others," says Dr. William Langston. But for many years, Parkinson's received special attention.
"For many years, the thought was Parkinson's was the perfect disease to lead the way in terms of solving these diseases," Langston explains. "The main reason for that is we were totally focused on one small area of the brain known as substantia nigra -- literally, black stuff -- [a] small pigmented dark area in the brain that sits atop the brain stem. Now, that looked like a pretty easy target, not a big nucleus. We fix that, we get more of the normal chemistry restored in the brain, and we fix the disease. To some degree, I think, that's still true. But we're now learning that Parkinson's is actually much more complex."
As scientists have learned that the symptoms of Parkinson's go far beyond the movement problems linked to the decay of the substantia nigra, they've realized that the simplicity that once made this seem like an easy neurological disorder to crack -- the best candidate for a cure -- was an illusion.
Does that mean what's learned about Parkinson's will have no effect on our ability to solve neurological diseases?
No. The complex nature of Parkinson's doesn't rule out the ability of breakthroughs studying this disease to have an impact on the understanding and treatment of other disorders.
"I think you're going to see most surgical therapies carried out first in Parkinson's disease," says Langston, citing the disease's pronounced effect on a small area of the brain. "And that's already happening with gene therapy, where genes are inserted in the brain to try to make cells healthier. ... If we get to the point of stem cells going in, all of that will probably be done with Parkinson's first."
Dr. Clive Svendsen is studying how to use stem cells to deliver growth factor -- a chemical that can regenerate important nerve cells in the brain -- as a treatment for Parkinson's. He points out that, "The Department of Defense has funded work in Parkinson's disease for a number of years. ... I think some of this comes through lobbying of people like Muhammad Ali and Michael J. Fox, and certain senators, to try and get funds appropriated specifically for Parkinson's disease. A number of veterans get Parkinson's disease. [With] Gulf War syndrome, there's an increased performance of ALS [another neurodegenerative disease]. And just neuro injury in wartime conditions is important to the Army. They're looking to stem cells releasing growth factors as a potential treatment for their troops on the field and for their veterans." Solving even a part of Parkinson's could still help solve parts of other brain disorders.
Are stem cells key to finding a cure?
Theoretically, stem cells should be able to replace the damaged brain cells whose degeneration leads to Parkinson's symptoms.
"We have shown very clearly that our basic science work in the laboratory ... proved that we can restore ... brain function in patients," says Dr. Ole Isacson of the Parkinson's Research Center at Harvard Medical School. "But we do need a lot of work to overcome the obstacles of making this happen [in] every case and in a reliable way."
"When we started with neurotransplants -- and we started with fetal cells first, and eventually the hope was stem cells would replace those -- we thought this was gonna be easy," Dr. William Langston explains. "We just put the cells in, fix this one small area of the brain, and we cure the disease. And we were very disappointed when that didn't happen. I think now with our evolving concept of Parkinson's disease, treating this one small area of the brain that we can already treat pretty well with surgical therapies, [stem cell therapy is] important but I think it is no longer the Holy Grail."
Some patients who've received experimental cell transplants have seen their symptoms improve, while others have experienced no change or even gotten worse.
"Some more extreme critics of this field," says Dr. Clive Svendsen, "would say that we now have two types of Parkinson's. We have Parkinson's disease and Parkinson's plus transplant disease."
[Editors' Note: Researchers at MIT's Whitehead Institute for Biomedical Research have developed a new method of reprogramming the skin cells of Parkinson's patients into "an embryonic-stem-cell-like state," using "the resulting cells to derive dopamine-producing neurons, the cell type that degenerates in Parkinson’s disease patients." According to The New York Times, this method, "would in principle allow the brain cells that are lost in Parkinson's to be replaced with cells that carried no risk of immune rejection."]
Why haven't stem cell transplants worked so far?
There are many reasons scientists have encountered obstacles in engineering successful stem cell treatments for Parkinson's.
"The comparison I always use is, imagine trying to wire your house after it was built," says Langston. "I mean, when you build a house, all the wiring goes in very early. When the house is built, if you had to do all the wiring afterwards, that would be pretty tricky. Now imagine you're trying to do that in a living brain with 4 billion neurons."
Some stem cells revert from the type of cells they've been cultivated to become -- generally dopamine neurons when dealing with Parkinson's -- back into generic stem cells when they're placed in the brain. And an adult brain seems to have ways of recognizing embryonic stem cells as not belonging. "We're learning there're all types of signals in the adult brain that tell these little guys to go away," says Langston.
Furthermore, a stem cell's inherent ability to develop into any kind of cell -- the property that makes them useful -- can also be a hazard. Says Isaacson, "While we know that we can generate the dopamine neuron, it also tends to generate other cell types, including skin, and maybe even bone. So the challenge is, even though we can get the cell we want, [we need] to eliminate the other unwanted cells, lest they would grow into tissues that would be very problematic and even dangerous to the patient."
In some cases, says Langston, "We don't even understand the reasons why the attempts we've already done have failed."
The risk of stem cell procedures compared to the safety of other relatively effective Parkinson's treatments gives some scientists pause. "The patients will have [to] think about this," says Svendsen. "Am I gonna risk a new procedure that hasn't been tried or go with the steady state procedures which we know work?"
Two perspectives on the promise of stem cells...
Because of the practical difficulties surrounding stem cell procedures, some scientists no longer think of them as a potential wonder cure-all.
"It's not popular to say that stem cells aren't the answer, but I now believe they are not," says Langston. "I don't think we should give up with them. I think they're gonna help other diseases, and eventually they may really help Parkinson's. But I don't think that should be our major focus. The brain is not a pincushion. We can't keep plucking cells in all over the place."
Others are more optimistic, although they point out it may take many years to fix the current problems with stem cell transplant technologies.
"What we're facing is the same kind of problems that you see almost in engineering. If you think of the early development of airplanes, or flight, the first airplane crashed very quickly," says Ole Isacson.
"But many years, in a way, is still a short time in my world, because if you think about the kind of discoveries that are necessary to make a medical treatment available to a large group of patients, decades is the norm rather than the exception."
"We're looking at this through the Parkinson's window," says Clive Svendsen, "but if you look outside the Parkinson's window and go to the rest of the world, [there are] lots of places that this technology will be used, I'm sure." |
Other Dairy Products
( Originally Published 1939 )
It has been known for a long time that babies fed with boiled milk thrived better than those fed with regular milk. It was supposed that this improved quality lay in the destruction of pathogenic organisms. Gradually a wealth of data has been accumulating which shows that boiling or other treatment of milk increases its digestibility. After Buckley had showed that the physical nature of the curd of milk is important in determining the food value of milk,' and Ladd had published some chemical data that showed that homogenized milk produced a soft curd in the infant's stomach and was similar to breast milk in this respect,' Washburn and Jones showed that homogenization of milk produced curds which were much more flocculent and friable than those of regular milk, although this property was not reflected in any improved nutrition of his experimental animals, young pigs. The recent work of Hill is credited with giving this subject of the digestive quality of milk an emphasis which has found important application in the commercial production of soft curd milk. Interest has been further stimulated competitively by reason of the inroads that the evaporated milk industry has made into the bottled trade, largely by reason of the superior properties of the canned product in infant feeding. This has led to laboratory activity directed toward devising processes for imparting soft curd properties and for measuring curd hardness rather than toward ascertaining to what extent, if any, these treatments actually improve the digestibility of the milk. The scientific literature leaves the subject in a very con-fused state. Whatever improved digestibility there is seems to result entirely from the speeding of the passage of the milk from the stomach and not from any increased food value or degree of assimilation.' The whole subject is excellently reviewed by Doan in the Journal of Dairy Science, 21, 739-756 (1938).
NATURAL SOFT-CURD MILK
Hill found that the milk of different cows possessed unequal digestibility, and that many infants could tolerate milk from certain cows but not that from others. In general, this improved tolerance was associated with milk of relatively low total solids content, al-though this relationship did not seem to be exclusively specific. Soft-curd milk was produced by cows of different breeds and was fairly uniform over the lactation period of a given cow. This property enabled herdsmen to select cows for the regular production of this kind of milk.
Soft-curd milk is more rapidly digested by humans, calves, and rats, and leaves the stomach more quickly than regular milks At the same time, soft-curd milk has a lower content of total solids and a smaller calorific value.
It has been observed that cows suffering with mastitis produce a soft-curd milk. This has led many persons to think that all soft-curd milk is pathologic. Such a belief is erroneous. Soft-curd milk is actually under more stringent control than regular milk because its production is mostly, if not entirely, limited to Grade A and certified herds. However, on account of the widespread prevalence of sub-clinical mastitis, it is recommended that the presence of udder infections be tested for when the curd tension is determined.
Elias- showed 10 that soft-curd milk gave curds in the stomach similar to those of boiled milk. Espe and Dye reported that doubling the curd tension increased the length of the digestive period from 30 to 65 percent, and that boiling markedly lowered the curd tension. Welch and Doan showed that curd tension was greater in milk of high casein content, and that equalization of casein content by dilution with water caused both the curd tension and differences in rates of digestion largely to disappear, although the casein content might exercise only a minor role in the rate of digestion if the curd is artificially softened by heating, homogenization, and other means.
ARTIFICIAL SOFT-CURD MILK
Soft-curd by homogenization. Softness of curd can be imparted to a milk by homogenizing it. This procedure consists in pumping milk under very great pressure through a special valve with small clearance so that the, butterfat globules are broken up and uniformly distributed. The homogenization of skimmed milk does not impart soft-curd properties; at least about 1 percent of butterfat or other oil must be present. Chocolate milk is a soft-curd milk. Therefore, it seems that the imparting of soft-curd properties by mechanical means is a function of the degree of dispersion of discrete particles whereby the curd is mechanically prevented from setting into a solid homogeneous mass. Feeding experiments on rats showed that this homogenized soft-curd milk was digested just about as quickly as boiled milk or natural soft-curd milk. Letters patent 12 have been issued to cover the production of soft-curd milk by homogenization, although the process seems to have been practiced by milk companies for many years previously to the granting of the patent.
The difficulty of controlling exactly the effectiveness of the homogenizing machine, together with the variability in the composition or physical nature of the milk, particularly the butterfat, precludes the determination of the most efficient temperatures and pressures. Experience has taught that the curd of a given milk cannot be softened beyond a certain point, regardless of the pressure used, and on the other hand, too light a pressure does not insure permanency to the imparted curd softness. In industrial practice, consistent results can be obtained when milk is homogenized at pressures of about 2500 to 3000 pounds per square inch at a temperature of about 145° F. This softens the curd to a tension of about 30 grams, or reduces the curd tension of average market milk about 50 percent.
The homogenization of milk must be carefully conducted if a satisfactory product is to be obtained. Trout and his associates found 13 that some milk upon homogenization developed rancidity within 15 minutes after treatment. This effect seemed to be caused by a lipolytic enzyme which could be inactivated at temperatures of pasteurization. Accordingly, this off-flavor can be prevented by pasteurizing the milk before or immediately after homogenization. The flavor of the finished product is generally considered to be slightly better if pasteurization precedes homogenization, but health officers are inclined to require pasteurization to come last.
Homogenized milk, unless the milk was initially of high quality, may exhibit a smudgy yellow or gray sediment in the bottom of the bottle. It is too finely divided to be revealed on a sediment disc. Babcock 14 reported that it consists largely of leucocytes, epithelial cells, and some finely divided dirt. Charles and Sommer 15 state that sediment may occur in milk of the highest sanitary quality and may come from a healthy udder. It is not seen in unhomogenized milk because the rising of the fat globules into the cream layer sweeps this light material upward. Clarification by centrifugal clarifying ma-chines will remove it.
Soft-curd properties, artificially imparted to milk by homogenization, were studied by Anthony on two adult males who possessed the unusual ability to regurgitate at will without distress. This enabled them to drink the milk, hold it in their stomachs for 30 minutes, and then return it without the aid of a stomach pump or an emetic. These experiments showed that the tests on curd strength made in vitro and determined with the curd knife reasonably evaluated the nature of the curd in the human stomach (except in the case of mineral modified milk). The curd particles of breast milk were minute and soft, and were so finely divided that they could not be separated from the accompanying juices with a 20-mesh screen. On cows' milk, when the curd tension (by laboratory curd-knife technic) was high, the regurgitation specimens of curd in every case were large and leathery. When the readings were low, the curd particles were small and soft. Breast milk registered 0 curd tension, natural cows' milk 50-100 grams, and homogenized milk (processed at 3500 pounds) 15 grams. The patients reported that the milk tasted better (because of the minute division of the milkfat globules) and gave less distress.
However, no digestive advantage is reported by some other investigators who worked on samples in vitro and on experimental animals. The latter work is not so impressive as clinical studies but may be better controlled. Much more fundamental and clinical re-search is necessary before the value of this processing is substantiated.
Soft curd by sonic vibration. A modification of the homogenizing process for the production of soft-curd milk has been developed by subjecting milk to intense sonic vibration. Electromagnetic oscillators, somewhat similar to those used in submarine communications, are constructed to allow the passage of milk in a thin film between the "anvil" and the vibrating diaphragm. Sonic vibration acts directly on the butterfat of the milk to cause a more complete dispersion. The reduction of curd tension is a function of the number of fat particles, and not of the actual fat concentration. A direct relationship seems to exist between the degree of fat dispersion-and the degree of curd-tension reduction. Inasmuch as only a small proportion of the total fat in milk need be finely subdivided to reduce the curd tension, it is possible to produce soft curd by vibration without destroying the cream volume (cream line).
Commercial homogenization. The practice of homogenizing market milk is gradually extending. It is quite general in parts of Canada, and is increasing irregularly in the United States. Fifteen states have no regulations for the control of homogenized milk, 19 states and the District of Columbia permit its sale if properly labeled, 2 states have taken no action but look upon it with disfavor, and 4 states prohibit its sale. It is a useful practice for the treatment of milk which is to be consumed in restaurants, institutions, or wherever the sale of bulk milk introduces the likelihood that the consumer may be served a portion from which a substantial part of the butterfat has separated. Tracy states that the unpopularity of homogenized milk in the past has been due largely to the emphasis placed on the cream line as a measure of the value of a milk, and to the unfriendly attitude of some regulatory officials who felt that homogenization might encourage fraudulent practices. About one-third of the milk-route customers of the University of Illinois changed to this milk for the following reasons: it looked and tasted better; no cream adhered to the bottle cap; no mixing was required; it tasted better for breakfast foods; it removed the temptation to abstract cream; it was easier to prepare for infant feeding; it did not allow rising of cream to top of glass in refrigerator; it made better milk drinks; it tasted better to children; it was more easily digested by infants; and it did not churn out on freezing.
Soft curd by base exchange. Hard waters are softened by the zeolite or base-exchange method, whereby the percolation of water through a bed of zeolite (a sodium-aluminum silicate) effects an ex-change of sodium and calcium. As applied to milk, sodium from the zeolite replaces soluble calcium from the milk. The milk is first acidified to about 0.3 percent as lactic acid (with a dilute nitric acid solution) and then percolated at 64° F. over a granular column of zeolite. During the process, the pH is adjusted to that of ordinary cows' milk (about 6.50), and the acidity is reduced to about 0.15 per-cent as lactic acid. This process is reported 22 to change the taste, appearance, and other qualities very little from those of regular milk. The cream line is practically the same as in pasteurized milk. Bacteria counts are said to be lowered by the filtering effect of the pass-age of the milk through the zeolite bed.
The Hill method cannot be used to measure the curd tension by this process because the Hill technic introduces about ten times as much soluble calcium into 100 milliliters of milk as is removed by the base-exchange treatment. Moreover, it is considered more desirable to use a method which more closely simulates gastric digestion. Such a method has been developed by Miller.
Hess, Poncher, and Woodward 24 studied the nutritional effects of such a milk on an infant on a metabolism frame. They report that, in spite of the decrease of the percentage content of total calcium and phosphorus, 100 milliliters of such milk per kilogram of body weight kept a normal growing infant in a positive calcium and phosphorus balance during the entire time of feeding.
Soft curd by enzymic action. Milk can be given soft-curd properties within a range of 20 to 30 grams by the addition of pancreatic extract, concentrated in the proportion of 1 part of the powder to 10,000 parts of milk. The milk containing the enzyme is heated at a temperature of 42° C. (108° F.) for 15 minutes, and then is pasteurized in the regular way. The preliminary heating brings about a partial digestion of the curd, and the pasteurization inactivates most of the enzyme. The mineral content, the protein, and the formol titration values remain substantially unchanged.
Standards of quality. The quality of curd is usually determined by the Hill test, or some modification of it. Although natural milks may give a range of reading on the scale from 15 to 200 grams of tension, the average of numerous milk supplies has been found to be about 60-70. The American Association of Medical Milk Commissions 27 specifies that a soft-curd milk must show a curd tension be-low 30 grams, determined at least twice at an interval of 1 to 5 days before it can be claimed to be a soft-curd milk, and that the test must be repeated at monthly intervals thereafter.
Determination of curd tension. Hill's method for determining the characteristics of milk curd is based on the measurement of the degree of toughness of the curd which is coagulated by pepsin in calcium chloride solution. The measurement is the indicated pull in grams necessary for a special knife to cut through the coagulated curd. The knife consists of several radial horizontal blades soldered at right angles to an upright slender rod. This knife is placed in a jar containing 100 milliliters of the milk to be tested. A coagulating solution of scale pepsin and calcium chloride is then added. This sets the curd around the knife. The knife is then hooked to a spring balance, and its pull as it cuts upward through the curd is read directly from the dial.
Caulfield and Riddell have shown that it is expedient to make each determination in triplicate, and that temperature of reaction and time interval between the addition of coagulant and cutting of the curd must be kept constant. Miller 23 has modified this method by substituting an acid pepsin solution for the pepsin-calcium chloride solution. The measurement of toughness of curd by this method substantially parallels the digestibility of the milk by animals. See also the method of the U. S. Department of Agriculture, and that of the American Dairy Science Association reported supra by Doan.
Determination of butterfat. Authorities are not in agreement as to the effect of homogenization on the accuracy of the butterfat de-termination by the Babcock method. Babcock found that in every case the homogenized milk averaged in fat about 0.1 percent lower than the same milk before it was homogenized. On the other hand, Tracy' states that homogenized milk can be tested satisfactorily by the Babcock method if both the acid and milk are at about 70° F., if the acid is added in small portions, if slightly less acid (1.5 milli-liters) is used, and if the solution is shaken well after each addition of acid.
Microbiological examination. Inasmuch as natural soft curd has been associated with mastitis, it is advisable in the interest of sanitation and wholesomeness to examine samples of natural soft-curd milk for the presence of mastitis organisms.
1. S. S. BucKLEY, Maryland Agr. Exp. Sta. Bul. 184, 1914.
2. M. LADD, Boston Med. and Surg. J., 173, 13 (1915).
3. R. M. WASHBURN and C. H. JONES, Vermont Agr. Exp. Bul. 195, 1916.
4. R. L. HILL, Utah Agr. Exp. Sta. Bul. 207, 1928; Circular 101, 1933.
5. Council on Foods, J. Am. Med. Assoc., 108, 2040, 2122 (1937).
6. F. J. DoAN and R. C. WELCH, Pennsylvania State College Agr. Exp. Sta. But. 312, 1934. See also F. J. DoAN and C. C. FLORA, ibid., 380, 1939.
7. H. C. HANSEN, D. R. THEOPHILUS, F. W. ATKESON, and E. M. GILDow,
J. Dairy Sci., 17, 257 (1934).
8. R. C. WELCH and F. J. DoAN, Milk Plant Monthly, 22 (11) 30 (1933).
9. W. V. HALVERSEN, V. A. CHERRINGTON, and H. C. HANSEN, J. Dairy Sci., 17, 281 (1934).
10. H. L. ELIAS, Am. J. Diseases Children, 44, 296 (1932).
11. D. L. EsPE and J. A. DYE, ibid., 43, 62 (1932).
12. R. FLUCKIOER, U. S. Patent 1,973,145, Sept. 11, 1934.
13. G. M. Tamil, C. P. HALLORAN, and I. GouLD, Mich. Agr. Exp. Sta. Tech. Bul. 145, 1935.
14. C. J. BARCOCK, U. S. Dept. Agr. Tech. Bul. 438, 1934.
15. D. A. CHARLES and H. H. SOMMER, Milk Plant Monthly, 24, 26, 32 (1935).
16. G. E. ANTHONY, The Bulletin (official publication of the Genesee County Medical Society), 9, March 4 (1936).
17. L. A. CHAMBERS, J. Dairy Sci., 19, 29 (1936).
18. Milk Dealer, 25, 36 (1936).
19. R. H. TRACY, Milk Plant Monthly, 24, 28 (1935).
20. U. S. Patent 1,954,769, assigned to M. & R. Dietetic Laboratories, Inc.
21. J. F. LYMAN, E. H. BROWNE, and H. E. OTTING, Ind. Eng. Chem., 25, 1297 (1933). Also see Milk Plant Monthly, January, 1934, p. 37.
22. H. E. OTTING and J. J. QuILLIGAN, Milk Dealer, 23, 36 (1934).
23. D. MILLER, J. Dairy Sci., 18, 259 (1935).
24. J. H. HESS, H. G. PONCHER, and H. WOODWARD, Am. J. Diseases Children, 48, 1058 (1934).
25. V. CONQUEST, A. W. TURNER, and H. J. REYNOLDS, J. Dairy Sci., 21, 361 (1938).
26. R. L. HILL, ibid., 6, 509 (1923).
27. Methods and Standards for the Production of Certified Milk, Am. Assoc. Med. Milk Commissions, New York, 1936.
28. W. J. CAULFIELD and W. H. RIDDELL, J. Dairy Sci., 17, 791 (1934).
29. Chief of Bureau of Dairy Industry, 1938, J. Milk Technol., 2, 48 (1939).
30. Curd Tension Committee. Rept. Annual Meeting Amer. Dairy Sci. Assoc.,
31. P. H. TRACY, Milk Dealer, 25, 30, 60 (1936). |
(Student produced study guide from Foss, Foss, and Trapp )
characterized by the use of the scientific method and based on the presumption that we (people) can know objectively and comprehend the objects around us. (146)
forms of understanding from which reality is to be deduced. (146)
refers to the philosophical movement and is not based on the problem of existence, but deals with the problem of words, and the ways of rhetorically thinking and speaking that were perfected as a way of philosophizing in the 15th Century. (147)
deals with scientific objectivity, universality, and rational deduction over other ways of knowing the world. (148) Cogito ergo sum: "I think; therefore I am." This is a self- evident axiom and is based on the human power to apprehend reality by means of reason. (146)
what in essence, separates the human being from the animal. (149)
ability to make adjustments in nature simply because we are humans. (149)
the basic process by which humans gain control over nature; refers to a basic capacity to grasp what is common or similar in ideas or experiences. (151)
"It leads to light because it stems from the need to see: that which is not obvious...is to be transferred," (154) and provides the connection between rhetorical and rational speech. (158)
fundamental or original principle upon which philosophical arguments are based. (155)
illuminates historical fact, making a situation concrete, relevant and understandable using metaphor and imagery. (157)
deductive in nature and achieving its effect through logical demonstration. (157)
a "system" of signs whose elements receive meaning through and within this system. Morse code is an example. (159)
superficial and mistaken definition of rhetoric, as a technical art of persuasion, that acts on emotions to form beliefs. (159)
practice, or doing. Reality is manifested in concrete situations.
coming to terms with things by studying words individually.
choosing what perspective to take in a situation; the unveiling of an essential meaning.
approaches in research which are concerned with human's role in constructing rhetorical knowledge rather than with the possibility of objective knowledge (165)
A. Grassi's education was the product of two opposing philosophic traditions: German Idealism and Italian Humanism.
1. Grassi's background of Italian Humanism was challenged at the University of Freiburg, where German philosophy dominated.
2. The dissonance of the two views led Grassi to examine his own beliefs more carefully, from which he determined that rhetoric constitutes the foundation of human thought.
B. Two people were especially influential in differentiating the two philosophies for Grassi.
1. Bertando Spaventa's (Italian philosopher of the late nineteenth century), following statement left an impression on Grassi. "The development of German thought is natural, free, and independent, in a word, it is critical. The development of Italian thought is unsteady, hindered, and dogmatic. This is the great difference." (p.145)
2. Martin Heidegger, a German philosopher who worked with Grassi for ten years in Freiburg, held a strongly negative attitude towards Italian philosophy. Heidegger's attitude was influential in causing Grassi to seriously consider the value of both German philosophy and Italian philosophy.
A. In order to understand Grassi's approach to rhetoric, more precise definitions are needed for the Scientific Tradition and Italian Humanism.
1. Scientific Tradition is based on objective knowledge.
a. Rational deduction is at the core of the scientific method and involves starting from the premises and deriving the inferences already inherent in them.
b. Grassi lists three limitations of this scientific paradigm, which he believes constrain what is studied as philosophy.
i. The scientific method examines first principles, but not their sources.
ii. The scientific method focuses on quantification.
iii Scientific thought is concerned only with universals.
2. Grassi's Italian Humanism refers to a philosophical movement.
a. Grassi's Humanism is Platonic and Aristotelian in orientation.
b. Grassi's Humanism is concerned with "the problem of words, metaphorical thought, and rhetorical thinking."
c. Grassi's Humanists sought to understand ways in which humans respond to a set of demands from the world and, by their linguistic choices, reveal the way they view this world around them.
B. The Scientific Approach was in direct opposition to Grassi's Humanist Approach. (p.147)
1. The Scientific approach deals with objectivity while the Humanist approach deals with distinctions and contextual variations.
2. The Scientific Approach came to dominate philosophy, while the Humanists were seen as searching for and moving toward this position.
3. Grassi believed that the Scientific Method was one tool for understanding, while Humanism dealt with broader areas and combined the areas of rhetoric and philosophy.
A. Grassi believed Vico represents the thought of Italian Humanism most fully.
1. Vico considered the rise of human history to be the basic problem of philosophy.
2. History is what differentiates humans from animals.
B. Grassi's support for Humanist thought is based in Vico's conception of the humanization of nature. (p.150)
1. Grassi has a term called "meeting the claims or demands of life."
a. All living beings experience the world using their senses, and inherently organize their environment to meet their basic needs.
b. Animals rely on instinct to function.
2. The human process is very complex.
a. Humans can choose and aren't limited to actions of instinct.
b. Humans can define images through language and therefore can interpret the world in different ways.
3. Humanization or historication of nature occurs when:
a. Humans become aware of these capabilities.
b. They begin to make adjustments in nature, or "direct their own destinies."
4. Humans must take sensory level meanings and translate them into an intellectual level.
5. The clearing of forests and the cultivating of land is the first unfolding of human consciousness. (p.151)
6. This feeling of control over nature wasn't a sudden change: there were three developing stages.
a. In the Cultural Age, humans felt they were a part of the cultural world.
b. In the Age of Heros, combination of heros and gods (superhuman benefactors) were seen as helping humans by introducing social institutions and laws.
c. In the Age of Humanity, humans realized that they could control nature.
C. Humans gain control over nature using the Ingenium, which is the process of humanization. (p.151)
1. Ingenium transfers meaning from the sensory world to a higher human one.
2. Ingenium frees humans by allowing them to see relationships and making connections in experience which are needed to think new thoughts.
D. There are three basic ways in which Ingenium is manifest to create the humans world. (p.152)
1. Imagination functions to grasp control of reality into two ways.
a. Imagination allows humans to realize that they are not bound to nature in the same way that animals are bound. (152)
b. Imagination allows humans to explain the world around them. It allows us to select certain interpretations of what we sense and allows us to define and order.
2. Work allows us to make and interpret connections of the sensory phenomena. Work allows us to act upon those interpretations made by our imagination. (p.153)
3. Language allows us to name and assign meanings to things in the world. By naming something we create a reality apart from the world.
E. Humanists sought to understand things in the context of practical human action.
1. Praxis is action: the application of abstract philosophical concepts into concrete situations. (p.153)
2. Grammarians examine words and interpret the abstract human condition in combination with individual action. (p.154)
A. Grassi refers to Aristotle and Cicero to define the metaphor.
1. Aristotle: "[The metaphor allows us] to see the similarity between what is actually the most widely separated."
2. Cicero's definition of metaphor said it was like a "light" which gives insight into a "relationship."
B. Metaphor transfers insight on several levels. (p.154)
1. At the most basic level, the metaphor allows us to grasp similarities between two unrelated things.
a. The metaphor operationalizes ingenium by allowing the human to connect himself/herself to the world of senses.
b. We relate to nature in human terms.
2. Language works metaphorically, transferring insight.
a. Language is symbolic because it helps us relate two dissimilar things.
b. Language helps us interpret and connect to our world and experiences.
3. The process of philosophizing is metaphoric.
a. A philosophical argument cannot be made without understanding the first principle.
b. First principles are nonrational and "experienced" as an "urge."
c. Philosophical systems are constructed with a first principle as the base.
d. The similarities we make between "urges" to understand philosophical problem and the actual logical arguments we use are metaphorical.
A. Grassi discusses the superiority of rhetorical language over rational speech.
1. Rhetorical language adapts various uses of imagery to illuminate historical fact and make it concrete, while rational speech is deductive and achieves effect through logical demonstration.
2. Rhetorical language deals with concrete particulars of life, while rational speech is universal and abstract.
3. Rhetorical language is like dialogue because it takes the world into account, while rational speech is monologic and has no need to interact.
4. Rhetorical language concentrates in images, symbols and metaphors, while rational speech is grounded in logical events and chronology.
5. Rhetorical language goes beyond a formal system, while rational speech is set in a "code" and can only move through the use of metaphor, which is indicative of rhetorical speech.
B. A third form of speech identified by Grassi is "external rhetorical speech."
1. This is the superficial and mistaken definition of rhetoric as a technical art of persuasion.
2. This is "false speech" because images do not stem directly from metaphors or nature, but a limited understanding of nature and its images.
A. Many feel that rhetoric is only the form of a message, while philosophy supplies factual content.
B. Humanists see rhetoric in a positive light, as a way to make logical reasoning palatable to an audience.
C. Grassi sees no separation of passion from logic.
1. The power of a message derives from its starting point in images that inspire wonder, admiration, engagement and passion.
2. Rhetoric, rather than logical deduction, is the true philosophy since it undertakes questions about the process by which "humans know, interpret, and create their world."
D. The emphasis on science in the Western world has resulted in this separation of content from form, and contrasts with the World View of the humanists. (p.161)
1. Without scientific proof, an idea will not be believed.
2. We have forgotten that we need to study the insights upon which these calculations are based.
E. There are many consequences for society that over-values the rational paradigm.
1. Those who believe in the "primacy of logic" and the ability for technology to deal with all problems tend to have an attitude of superiority.
a. Humans see their rationality as giving them a dominance over all things.
b. This actually limits humans' capability to fully interpret all things.
2. This affects our relationships with other cultures who do not share this attitude.
a. We see these cultures as being underdeveloped.
b. This view makes it impossible to fully understand them and constricts our interactions with them.
3. Logical thought has become synonymous with the domination of humans. (p.162)
4. The dawning of the atomic age is the ultimate example of humans' need to dominate nature.
F. The rational approach which has dominated Western culture has been detrimental to philosophy.
A. Renaissance Humanism defined folly as speaking irrationally without reason.
B. Grassi studied literature for examples of folly, defining it instead as the ability, using language, to choose the perspective to take on a situation to unveil something's essence. (p.163)
1. Folly is an engagement of ingenium
2. Folly is the fundamental process by which humans move from the nonhuman to the human realm.
3. Folly, as an extension of ingenium, allows humans to imagine themselves in new situation and to deal with these situations effectively.
A. Grassi's ideas of rhetoric are not well known among communication scholars.
1. Grassi has published in English only in Philosophy and Rhetoric and there are few essays or discussions on his work.
2. Grassi asserts the contributions of Italian Humanism to rhetoric and philosophy rather than with fully developing the contemporary implications of the philosophic perspective.
B. Grassi made several important contributions to rhetoric from the Humanistic perspective.
1. Grassi asserts that rhetoric and philosophy are necessarily connected, since rhetoric is the starting point of philosophy.
2. Grassi preference of thought, speech, and action made from connections with nature (ingenium) rather than from logical reasoning is similar to the "new paradigm."
3. Grassi's notion of folly allows humans choices in how they perceive the world they live in.
4. Grassi's work generated renewed interest in Renaissance Humanism.
5. Grassi gives new significance to rhetorical speech and asks us to reconceptualize our definition of rhetoric.
Grassi, Ernesto "Italian Humanism and Heidegger's Thesis of the End of Philosophy," Philosophy and Rhetoric, 13 (Spring 1980), 83. In this article, Grassi points to the parallelism between Heidegger's German Idealistic thought and the Italian Humanist tradition in order to create a historical framework in which to make evident the problems of Humanism in relation to present day. This article, separated into ten major ideas, begins with the End of Metaphysics and ends with Heidegger's Theory of the Brutality of the Being. In between these two major points, as Grassi explains Heidegger's Twofold thesis, the Traditional Model of Scientific Thought, the basic problem of Italian Humanism, The Question of the Veil of the Poetic Word, and the "clearing" of the Primordial Forest.
Grassi, Ernesto Die Macht des Bildes, 221, cited in Walter Veit, "The Potency of Imagery - the Impotence of Rational Language: Ernesto Grassi's Contribution to Modern Epistemology," Philosophy and Rhetoric, 17 (1984), 235. Veit gives analysis to some of Grassi's theories in this article. Grassi once again confronts the separation of logical reasoning and rhetoric. Many of Grassi's ideas have literally reconstructed the philosophical dimension of rhetoric in the eyes of contemporary Italian Humanists. Much of this article is similar to the subjects covered in FFT, and it even helps in the understanding since it gives different explanations and examples to similar material. Emphasis is placed on some of the ideals of eighteenth century philosopher Giambattista Vico, who Grassi found as a source for some of his rhetorical ideas. Once again, it is stated that Grassi believed that the philosophical revolution began with the Italian Humanists, who showed that philosophy gains insight into the principles "through the creativity of the image."
Grassi, Ernesto "Humanistic Rhetorical Philosophizing: Giovanni Pontano's Theory of the Unity of Poetry, Rhetoric, and History," Philosophy and Rhetoric, 17 (1984), 146. This article is simply Grassi's analysis and reaction to Potano's theories. It gives the reader some idea of the process of critical thinking that Grassi goes through in regards to the ideas of others. It is standard to differentiate between logic and rhetoric. The premises resulting from a rational process as exemplified by traditional metaphysics are necessary and universally valued. Rhetoric is bound by time and place, and it must use metaphor and images in order to be effective. In order for metaphor to be effective, there must be a common viewpoint shared between source and receiver which permits the audience to see the relationship of the metaphor. The unity of poetry, rhetoric, and history has a philosophical significance. All three are rooted in directive language. Potano's ideas are that the traditional thoughts need to be revised. A new kind of philosophy starts with the Humanists and the turn to rhetoric, away from rational argument.
Grassi, Ernesto "Remarks on German Idealism, Humanism and the Philosophical Function of Rhetoric," Philosophy and Rhetoric, 19 (1986), 125. Grassi discusses his blending of Vico's Italian Humanism and German Idealism. The entire system of thoughts is summarized as follows: "The faculty that is crucial to the making of metaphors is ingenium, which allows us to see the world. The power of language is beyond logic and rational thinking. To think rationally involves assuming some presuppositions and drawing inferences from them." This text is simply and expansion of the work in Foss, et al. and provides a historical context for Grassi's work.
Grassi, Ernesto "The Ordinary Quality of the Poetic and Rhetorical Word: Heidegger, Ungaretti, and Neruda," Philosophy and Rhetoric, 20 (1987), 248. This article is divided into three sections. The first section concentrates on making clear the philosophical function of poetical and rhetorical language, by looking at statements of philosopher Martin Heidegger and of two poets Ungaretti and Neruda. The poetic world, according to Heidegger, receives not only priority over the rational world, but also has a philosophical function comparing it to the ideas of philosopher Giambattista Vico. Second, it deals with the idea that reality cannot be revealed through a rational process. Next it deals with the philosophical function of poetry, showing that every beginning of a historical era is announced with a poetic expression, showing connection of poetry, rhetoric and history. Imagery is poetry. The third part shows what can happen when rational word becomes superior to rhetorical word, using the stories of Prometheus and Ulyssses. The fire Prometheus brought is considered metaphoric, but the fact that his liver is being destroyed keeps him historical and not eternal. Dante condemns Ulysses when he wishes to go beyond Hercules' pillars.
Grassi, Ernesto "Why Rhetoric is Philosophy," Philosophy and Rhetoric, 20 (1987), 75. Traditional philosophy arrives at an important admission: rational language cannot reach "passions." What is "true" language? The model provided by German romantic thought recognizes an essentially literary character. In Monologne, Noralis, language is a game; language does not occur for the determination of beings. Tongue speaks for itself alone. An object has its own destiny and at the same time it doesn't in that each appears in its merry through the code which is revealed in the history. Rhetorical, historical language is shown to be the true philosophical language because it is by means of it that we "uncover" the various world by "playing" with our "orders" at stake.
Verene, Donald Phillip, rev. of Die Macht der Phantasie and Rhetoric
as Philosophy, by Ernesto Grassi, Philosophy and Rhetoric, 13 (Fall
1980), 281. This article summarizes Grassi's ideas in Rhetoric as Philosophy:
The Humanist Tradition and in Die Macht der Phantasie. Zur Geschitchte
abenlandischen Denkens. According to Verene, these books are "treasure
houses of an understanding of the nature of rhetoric and its relationship
to philosophy that is absent in contemporary thought." Grassi's thesis
is that rhetoric is at the basis of philosophy. Considering this relationship,
Grassi asks his readers to understand the power of
language by choosing Humanism over science.
Verene, Donald Phillip, "Response
to Grassi," Philosophy
and Rhetoric, 19 (1986), 135. Verene is delighted to be discussing Grassi's
work. He admits skepticism to the blending of Italian Humanism and German
Idealism. The most important element of Grassi's thought is the metaphor
because metaphors embody the starting pints for thought. It is essentially
a recovery of ancient ideas that Verene feels is long overdue.
back to lecture note index |
Commercially-pressed CD's and CD-R or CD-RW disks are fundamentally different technologies, which is why a commercial CD will continue to be readable long after a CD-R has become unusable.
A CD drive uses a focused laser beam that is reflected from the media surface in the CD disc. The beam is reflected onto a sensor that detects changes in the amount of energy that is reflected. The original (commercial) process used perforated aluminum as the media surface. When you use the term "pressed" you are using an old vinyl record term, but the production process is pretty much the same. There is a "master" disk that is put into a press which is filled with polycarbonate. The master disk has little pins sticking up everywhere there is to be a hole in the aluminum. The disk is cooled, and liquid aluminum is spun onto it. This results in an aluminum layer with holes in it.
When the disc is played, the laser reflects strongly from the shiny aluminum or less strongly (or not at all) from the holes. The reflection/non-reflection is translated into the ones and zeros of the binary data stored on the disk.
Over time the aluminum can oxidize or there can be other changes in the plastic and other materials that make the disc unusable. These are long term effects and the ultimate statistical life of a commercial CD is often debated, without conclusion, by the experts.
The CD-R and CD-RW do not use an aluminum media surface. Instead, they use a dye. When the disc is written, a high-powered laser causes spots on the disk to turn dark (hence the term "burning"). When played back, the sensor in the player sees the difference in reflectivity of the dark and not-so-dark spots as the binary data.
Unfortunately, because the dye is a light-sensitive chemical, over time it will fade. This can happen from the heat of the reading laser, from ambient light, and from chemical degradation in the dye and support media.
CD-R/RW media is safe for backup, and for creating alternate media (copying music files to play in your car so that if they are damaged from heat or wear out you can make another one, and preserve your originals elsewhere), and similar purposes. However, they are not safe for archival storage because they are not stable enough for that purpose.
Side note: when burning CD's for use in a car, for best results get "music CD's" which are designed for that application, or slow your burning speed down to 12x or 16x to get a darker spot from your high speed burner. The car will read the disc more reliably.
Insofar as tape storage is concerned, tape is also not a good archival choice of media. It's generally better than CD-R, although I haven't seen any comparative studies.
Major data centers who use tape storage refresh the storage periodically. Their Tape Management System (TMS) remembers the date the tape was recorded, and will call it up to be copied periodically. The old tape is then erased and reused until it reaches end of life (sometimes a fixed usage or time interval, sometimes when the number of recoverable errors reaches a threshold) at which time it is scrapped.
The whole issue of long term archive is complex, and goes beyond media. For media, if a data center stored its files on a 9 track magnetic tape twenty years ago, how would it retrieve that data today (you cannot find working 9 track drives). What if it had used an early Magneto-Optical (MO) drive? Small businesses have trouble when their tape drive fails, and they can't buy another drive in that old format.
File formats are another problem. I have word processing documents that deceased family members created years ago. I no longer have word processing software that will import some of those formats. I can (sometimes) extract the raw text and then try to reformat it in a current program, but if I don't have a printed original I don't know how it was intended to be formatted.
The only archival format that has stood the test of time is paper.
Submitted by: Kevin G. of Dallas, TX
Well, Carl, that so-called expert sure has stirred the waters and a LOT of people are wondering about the same question. However, your friendly Federal Government has studied the problem even longer.
To be specific, the National Archives and Records Administration, in charge of all of the record archiving of the government, has no standard on media storage, and requested NIST, that's National Institute of Standards and Technology, to write a new standard on media durability.
If you never heard of NIST, you're not alone, as NIST is more of a background organization, but suffice to say, they're the ones who creates the standards, references, and accuracy tests for all industries, from DNA to Time accuracy (in fact, if NIST operates one of the Internet "clocks" you can calibrate your PC to). NIST DNA reference material improves forensic DNA test accuracy. NIST also invented closed captioning and many other technology, but enough about NIST.
A gentleman by the name of Fred Byers spent a whole year testing various media, and wrote a guide for NIST to librarians who need to archive information on how to care for optical media such as CD-R and DVD-R's and such. In the guide, he basically stated that with proper handling (store in low humidity, no scratching, stored vertically, etc.) a DVD-R should last 30 years with no fear of losing any information. However, that is NOT an absolute number as it is dependent on a LARGE NUMBER OF FACTORS, some of which in your control, and some not:
Factors that affect disc life expectancy include the following:
type -- as recordable media is more durable than rewritable media manufacturing quality -- you get what you pay for condition of the disc before recording -- obvious quality of the disc recording -- garbage in, garbage out handling and maintenance -- scratches are bad for any discs environmental conditions -- humidity and temperature can warp disc, ruining the reflective layer in the media. light, esp. UV light can destroy the dye used in recordable media, etc.
Let us discuss each factor in a bit more detail
All types of media can be damaged through warpage (disc bending), scratches, and reflective layer breakdown due to oxygen leakage.
Recordable media, in addition, is susceptible to UV rays, which affects the dye used in the process.
Rewritable media, with phase-change recording, is even more susceptible to UV ray and temperature.
It is generally acknowledged that certain brands of media are better than others, and often the stuff on sale is not the stuff you may want to buy and keep around.
What you may not know is that there are only like 16 media manufacturers in the world. They make the media for all the brands that you see in the market, and some brands / factories are known to make high grade media (i.e. they tested best for maintaining data integrity, even when the media was subject to aging tests). While few independent labs did comprehensives tests, a test in Europe a while back for CD-R's revealed that Taiyo Yuden (Verbatim), Kodak (Kodak), and TDK (TDK) kept the most data intact.
Condition of the disc before recording
A disc should be brand new when used. While shelf live of a media is up to 5 years, why take chances? Buy them as you need tem.
Handling and maintenance
Scratches are bad for any discs, as it breaks open the substrate layer and allows air to tarnish the inside silver reflective layer inside.
Scratches also can make information on the media unreadable by interrupting the laser's path.
Environmental conditions -- humidity and temperature can warp the media, and exposure to UV light can destroy the dye used in CD-R's and DVD-R's.
Hope that answers your questions.
Submitted by: Kasey C. of San Francisco, CA
In the '80's, the CD was introduced in the market and portrayed as "THE" solution to the vinyl records.
The CD could be thrown in a mud pool, step on it, scratch it, nothing would harm he CD.
Now we all know that CD's has a lower lifetime as their vinyl counterparts and are more susceptible to errors than them. This is also true for the CD as a media to record software.
The early CD's, were recorded at maximum 640 MB. Mostly not even 640 MB but something like 528 MB. This made them less susceptible to scratches.
But as the CD technology was in a constant evolution, overburning a CD to 800 MB and more became common use. Also the DVD was introduced, offering 4/9GB on a wafer of the size of a CD.
It is obvious the the tracks are becoming so small that the finest scratch, the smallest fault, can ruin the CD/DVD forever.
Answering your question, there is no miracle solution to keep CD/DVD from deterioration trough age. But with a little bit of care, you can have many years of pleasure of your recordings.
1. Buy only CD/DVD from a good brand.
Buying low priced CD/DVD will mostly result in very disappointing experiences.
2. Dont overburn a CD/DVD.
While the overburn technique is now widely accepted by most software, it is still not fully reliable and mostly dont approved by the CD manufacturers.
3. Put every CD back in the jewel case after use, clean them as prescribed by the manufacturer, and avoid as much as possible touching the reading surface of the CD.
As a final remark, CD/DVD are nowadays not expensive and if you can make a backup of them, make a backup and store it in a safe place.
I use an external harddisk of 250GB (<300 CD's) to store a backup of the CD's/DVD's I have. Price of the harddisk is about 100$.
I have been able many times to rescue recordings which were otherwise lost forever by this harddisk.
Hope this helps,
Submitted by: Carlos
You have just discovered what most people don't discover until they actually lose data: commercial CDs and home-burned CDs are not the same. While a commercial premade CD will last a very long time if it is cared for properly, a home-burned CD will begin to deteriorate. The reason is that the home-burned version uses dyes to accomplish what the premade CD does by having it built into the disk. This is, of course, an oversimplified explanation, but it will suffice.
There are a few ways to maximize the amount of time a CD will last. First of all, buy good quality CDs to begin with. Stick with brand names that you are familiar with and have used successfully in the past.
Do not assume that just because a blank CD is made by a well-known company that it will be high quality.
Test them out by actually using them. One of the best ways to do this is to use them for your regular system backups. Be sure to actually restore from those backups periodically (easier if you have another computer handy that you can wipe out data on) or else use a backup program that allows you to mount the backup as a "virtual drive" and retrieve data from it.
This lets you know if there is a problem with a brand deteriorating unusually fast.
Second, never use labels on CDs. I found this out the hard way. Labels cause the CD to deteriorate much more rapidly than it otherwise would. Certain inks used in pens have been reported to do the same, but I have never encountered this problem, so it shouldn't be too severe. Do be certain, however, that you are gentle when marking CDs. Use a felt tip and do not press hard.
Third, put a note somewhere on the CD that tells you when you made it. This lets you monitor how long it has been since the CD was burned. If the data is irreplaceable, burn it to a new CD every 2 years.
As for the recommendation to use magnetic tapes, that has its own set of problems. Magnetic tapes also deteriorate, and they are subject to some damage that CDs are immune to, notably damage from electrical or magnetic fields.
In short, CDs are good for long term storage-- but don't assume that "long term" means forever. Check them regularly and burn them to new media when problems develop or even before if you can't replace what's on them.
As for storage, that is pretty much common sense. Keep the CDs in a case or an envelope if they are not actually being used. Avoid temperature extremes and handle gently. I also recommend making two copies of every important CD. This practice just saved my data when I discovered that the labels on my CDs had wiped out some irreplaceable family photos. It costs twice as much, but if the data is important to you then it isn't really very expensive, is it?
Submitted by: Denise R. of Lebanon, Missouri
Hi Carl N,
Your question has been set by a lot of people over the last 10 years. I've burned CD in 100s over the years and only found 2 discs with missing information. Lifetime of CDs is not limited due to one parameters only, more issues are setting the limit of lifetime. One is related conditions of storage and how you handle the discs. In other word, how careless or careful you are as the user. Then the material used in the CDs - how cheap a blank CD did you buy. And lastly your burning equipment, that is the laser diode.
When pressed CDs were introduced in beginning of the 80s lovers of vinyl records claimed, that CDs would last for 2 years only. But as you have experienced CDs from this period can still be played. I remember one report from about 1990, which claimed a lifetime of only 3-4 years. Looking into the report, it turned out that the condition of storage was -30C (some -25F) and reading/playing equipment had a worn laser diode. Most of us can only say: I don't store my CDs in the freezer and todays laser diodes doesn't wear out as they used to.
Turning towards recordable CDs, the whole issue is a matter of having a whole bunch of clear holes placed in circles in a foil. Readability of these holes are depending on a number of issues. How clear are the holes? Is edge of the hole clear? Is the reflectivity of the materials sufficient? Is the laser still as effective as it was? or has the surface become matte? For the early CD-burners this was jeopardized by increasing burning speed and some blank CD had doubtful foil material. Adding, to this some CD-burners were even sold with writing speeds beyond its capability. Many blank CDs were rejected in this period due to bad burners rather than bad discs. It is my conclusion, that this interim period has given us some doubtful discs.
You have to be careful with your CDs and also a little bit extra careful with your own. They don't like heat, bright light, bending, and writing with aggressive writing pens is also nasty. Especially pens with unknown chemicals may etch the CDs, it is just like burning, but this time controlled by the chemicals.
In case you wanna increase the lifetime of your recordings, you may buy CDs which is claimed to 300+ years lifetime. These CDs are referred to as GOLD-CD. These CDs has a special layer which include some 24ct gold. The advantages of these CDs are the ability to create clear holes into it with reduced oxidation or corrosion over time. Amazing almost also unbelievable 300 year. Just 100 years could be great for me. In 10-20 years everything would be transferred to new media type anyway. I saw a report on the 300 years at
Price of these GOLD CDs is 10-20x times the usual ones.
You have been suggested to use magnetic tape. Nor tapes does not last forever. As matter of fact the sound quality decays over time; frequency range is decreased by each use. This loss you can not be restored as with a digital media. Only digital storage keeps its audio frequency range over time and use. Like R&R;DIGITAL media is here to stay. You may call it CD, DVD, MPx, Blu-ray or whatever, but it's digital.
I believe that todays discs and equipment can provide a disc with sufficient lifetime for most of us and may even restore your more doubtful discs from the early burning time with success. Even discs which are registered as 'No disc' may be restored by copying it today. In case you wanna assure yourselves; let the PC verify the burned disc, this option is normally disabled by default.
What shall I do with my precious discs from the early days? My best recommendation would be to make a new copy, while the old is still readable. This is easy and cheap to most of us today as having two drives in your PC is not uncommon. Lastly, the quality and lifetime of recorded discs is today likely to most depended on your own care.
Submitted by: Leif M. of Helsingor, Denmark
Regarding problems developing over time with recordable CD-R media, I've run into some of this myself, but I also have quite a few discs that were made back when the very first 1x CD burners were made available to the public, and they still read just fine for me.
I suspect that there are several factors involved here.
1. I'm certain there's a difference in quality between brands of CD-R media. A number of my really old CD-Rs that still read flawlessly today were Kodak branded, and were considered expensive "premium quality" discs at the time. They're even physically a little bit thicker than most other media I've handled. By contrast, some of the generic media I purchased because of the low price on 100-pack spindles has actually developed "bubbles" where you can see the dye that's sandwiched between the layers of plastic is disintegrating. (Of course it won't read if small spots are completely gone!) There were/are several different types of dye used for CD-R media, as well, and I wouldn't be surprised if it's turning out that some types have better longevity than others. For example, Verbatim was known for using their trademark blue-tinted dye, while others were shades of green or gold.
2. From what I've read and observed, handling makes a big difference too. Leaving your CD-R's exposed to sunlight (as folks tend to do with music CDs used in their cars or trucks) probably shaves years off of their lifespan. Putting them in some type of jewel-case or sleeve when not in use is a very good idea. Boxes of empty jewel-cases can be purchased fairly inexpensively at most office supply and electronics chain stores.
3. A CD-R holding computer data is inherently more "fragile" and subject to data loss than a CD-R recorded as an audio disc. The standard used for recording audio CDs incorporates quite a bit of error correction information to handle small scuffs and scratches on the media, but besides that, audio data is spread out over a much larger portion of the CD-R. If you have a .ZIP file stored on a CD-R, for example, a pinhole-sized mark someplace on the disc where that .ZIP file is stored can easily be enough to prevent the whole archive from extracting properly. By contrast, the same sized mark might only cause a very brief "stutter" at one point of a song on a music disc (or not pose a discernable problem at all, due to the error correction).
If your audio discs are already deteriorating to the point where players are rejecting them as "unreadable" or they're skipping badly, it sounds to me like things have gotten pretty bad. The only recommendation I'd have is to re-record your music to fresh, good-quality CD-R media and throw out the old ones - and in the future, make a habit of transferring your music to fresh discs every few years or so.
Luckily, in the case of computer software backups, they tend to become so outdated, you no longer really need to keep them by the time the media they're recorded on starts failing. But for those trying to preserve digital photos and the like, I'd recommend this same procedure. Make a fresh set of backups every so often and discard the old media - before it fails on you and you lose something priceless!
Submitted by: Tom W.
CD burned media fails after time.
I am a practicing technician and this is not a new complaint. It is my firm belief that most consumers burn their media at the fastest speed possible for both their software and the media they use. This is fine but there may well be a trade-off in doing this.
What most consumers do not perhaps understand is that commercially produced CD's have actual pits pressed into them that represent the digital data of the original sound data. A burned CD on the other hand is made by fabricating a photo sensitive layer to mimic the pits found in pressed media.
I have found three major causes for this consumers problem they are as
1. A slower burn makes a stronger image representation in the photo
sensitive layer of a burned CD. A faster burn while successful may not impress the photo sensitive layer as effectively as a slow burn. Over time the burn fails as the photo sensitive layer deteriorates.
2. Sunlight and other forms of intense light can effect a burned CD
because it can cause a distortion in the burned media's photo sensitive layer.
3. Scratches by far are more evident on burned media and more easily
caused than on pressed media. Most consumers seem to ignore the manufactures warning and suggestions. Handling of the disc in a careful manner as advised by the manufacturer is the best policy here. I use a camera lens cloth to clean surface of all my media. A camera lens cloth will not scratch the disc surface. Paper and regular household cloths will cause scratches.
Observe the above and I do believe you will have better results.
One more thing always use the media recommended by the burner manufacturer. It is endorsed and guaranteed to work, many of the cheap non name discs out there are just not up to par. Its just like the old cassette tape days.
Most audiophiles went for tapes like Maxell, JVC, Sony etc. but as everyone knows there were a lot of bogus brands out there for the un-informed to purchase.
Submitted by: Peter K.
> I recently read an article by a data storage expert who claimed that
> burned CD-Rs and CD-RWs can be expected to last only two to five years
> and not a whole lot more. I personally have commercially pressed CDs
> from the 1980s that still play fine, but I have begun to notice that
> some of my burned CD-Rs are beginning to skip
you mention that there are basically two types of CDs: Those that are created with all information in place and those you can buy and write on.
The first type is quite robust as the information has been "engraved" into the surface just below the reflector. The most critical part of such a CD is the reflector, most often a very thin layer of aluminum.
The second type of CD works a bit differently: There is a dye layer below the reflector and the information is written onto the CD-R(W) by "burning"
and thereby locally changing the optical properties of the dye. The most critical part is the dye, besides the reflector as above. If the dye degrades the CD easily gets unreadable. The dye of CD-RW is even more critical as it must be "resetable" - another constraint.
> The expert suggests that for secure long-term storage, high -quality
> magnetic tape is the way to go.
This solution is quite expensive as you need a tape drive and enough tape cartridges, but has the advantage of a much larger storage capacity. If the manufacturers say their tape cartridges are reliable for a very long time they have one advantage above CD-R: This type of storage device has been around long enough to prove it. CD-R has been on the market for no more than 10 years.
The best strategy for the private user is: Have a good archive strategy, save often and store the media carefully in a dry, dark, cool place. If you store every file more than once you have a better chance to retrieve it.
There is no real alternatives to CD-R. Use high-quality ones. Do not use any DVD variety as their reliability is much less. DVD may be used for an image backup of your boot drive so you can restore your present configuration for the months to come.
Submitted by: Alexander V.
Unlike pressed original CDs, burned CDs have a relatively short life span of between two to five years, depending on the quality of the CD. There are a few things you can do to extend the life of a burned CD, like keeping the disc in a cool, dark space, but not a whole lot more.
The problem is material degradation. Optical discs commonly used for burning, such as CD-R and CD-RW, have a recording surface consisting of a layer of dye that can be modified by heat to store data. The degradation process can result in the data "shifting" on the surface and thus becoming unreadable to the laser beam.
Many of the cheap burnable CDs available at discount stores have a life span of around two years, In fact, there are some of the better-quality discs offer a longer life span, of a maximum of five years. Distinguishing high-quality burnable CDs from low-quality discs is difficult, I think because few vendors use life span as a selling point.
I've had good luck with Verbatim media, and bad results with TDK. Playback with the TDK discs I used degraded steadily over time, in spite of very little use, and not much in the way of scratches or other blemishes on the disc. On the other hand, the Verbatim discs I've used have held up well over time, and under more use than the TDK ones I used.
Opinions vary on how to preserve data on digital storage media, such as optical CDs and DVDs. I have my own view: To overcome the preservation limitations of burnable CDs, Im suggesting using magnetic tapes, which, as I read, can have a life span of 30 years to 100 years, depending on their quality. Even if magnetic tapes are also subject to degradation, they're still the superior storage media.
But I want to point out that no storage medium lasts forever and, consequently, consumers and business alike need to have a migration plan to new storage technologies.
A Good Question to get in this subject is Does Burning Speed makes difference in quality of CDs? Someone told me that the burning speed makes a difference in the quality of the records. The lower it is, the deeper it burns and therefore the better the quality is. I heard that there are some audio technicians decide to burn masters at 2x and copies at 4x due to getting digital noise from higher burn rates. Might just depend on the burner quality and the burning program...
Hoping you get the Point of my explanation.
Submitted by: Sameer T.
Everyone who owns a business is always trying to be enticed by the security and the longevity of magnetic tape. And although I'm apt to agree with them on its durability, I don't use it to back up important data in my business. I have two problems with magnetic tape vs. CD or DVD. The first problem is hardware. Data backed up on a CD or DVD can be loaded into any computer with a drive capable of handling the media. The same can be said about tape backup, however you are more likely to find a computer off the shelf with a compatible CD or DVD drive vs a magnetic tape drive. The other problem is the need for long term storage.
As a business owner, I'm backing up my important data every one to three days. I've been using CD-RW media to do this for years. If a disk gets corrupt, you can reformat it using your burning software, then use it again. If you are concerned about your CD becoming corrupt, simply burn two or three. The cost of three CD or even DVD media is much more reasonable than the cost of one of those tapes. And if my server dies, I can buy any computer at any store, and load the data onto the new computer right away and I'm back in the game. I have several people trying to convince me of the benefit of a paper backup system. I find it easier (and cheaper) to have multiple electronic backups. My business server has a RAID 1 card and two hard drives which mirror each other. I have the CD backup, and I then take this data and save it to a secure partition on another machine in a separate location.
The likelihood of all of these systems failing at one time is highly unlikely. And if it does, I'm taking the day off, because that's just real bad luck. As far as long term storage (like music), I've noticed that CD-RW can lose the data in the long haul, but haven't had any problems with CD-R media. I have some that skip, but not for a reason I didn't know about. I buy the large spools of CD-R which don't come with jewel cases. So these disks get abused. If you know you are going to keep something for a really long time, I would make sure the disks you buy have jewel cases. And you can apply the multiple disk system with this as well. The media is easy on the wallet, and the more backups you have, the lower the risk you will actually lose the data.
I don't think I answered all the questions, but that's my take.
Submitted by: Dave K. |
Biology 107 Lab Exam Review
Science B01 with Magor at University of Alberta
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Size: 132 flashcards
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When is a plant cell plasmolyzed?
What is a hypertonic solution?
A hypertonic solution is a solution that has a higher concentration of solutes than found inside the cell
What are some of the functions of cell membranes?
Separation of cell contents from external environment
Organization of chemicals and reactions into specific organelles within the cell
Regulation of the transport of certain molecules into and out of the cell and its organelles
Beet cells red pigment, located in the cell's large central vacuole and surrounded by the tonoplast membrane
Molecule or part of a molecule that absorbs radiant energy (light)
Graph that shows the amount of light absorbed at a number of wavelengths
Four parts of a spectrophotometer
Device that isolates a photoelectric tube
Standard curve limitations
Standard curve is specific to the pigment and its buffer
Standard curve cannot be used for absorbances beyond the range of the standard curve
Cu = Cd/D
Deposits single bacterial cells from a liquid culture over the surface of an agar medium
Prevents bacteria in the environment from contaminating work, and prevents bacteria in work from contaminating the environment
Aspetic techniques that can be used
Sterilize surfaces and working surfaces
Washing hand before and after
Flaming inoculating loop and lip of culture tubes
Reduce time sterile medium, cultures, or bacteria are exposed to air
Work in area with low resident population of bacteria
What does the blank used with a spectrophotometer consist of?
How is a spectrophotometer zeroed?
Define each term in the "fluid mosaic model"
Fluid: membranes are able to move, things may pass through them (selective permeability)
Mosaic: membranes are composed of a variety of different things - phospholipid bilayer, enzymes, proteins that act as channels
Explain phospholipid bilayers
Fatty acid bilayers, where the hydrophobic ends of the fatty acids attract each other to the inside of the layer and the hydrophilic ends are on the outside of the membrane in the water, creating a double layer of fatty acids
How are bacterial species identified?
Cell and colony morphology, chemical composition of cell walls, biochemical activities, and nutritional requirements
What is the best way to isolate individual cells?
Streaking them onto an agar plate, so that each individual cell will produce a colony
Cell wall is composed of a thick layer of peptidoglycan surrounding the cell membrane
Differential stain to divide bacteria into Gram-positive and Gram-negatie
Steps of a gram stain
A basic dye (crystal violet) is used to stain the peptidoglycan in both cells, then iodine is used to increase affinity of the dye to peptidoglycan. Ethanol is then used to dissolve lipids in outer membrane of Gram-negative bacteria, allowing iodine dye complex to leave cells (peptidoglycan layer too thin to retain dye), while Gram-positive cells retain the dyed due to the thick layer of peptidoglycan - a counterstain is then applied that dies the Gram-negative cells pink
Commonly reocognized cell morphologies
Cocci: spherical shape
Bacilli: shaped like rods or cylinders (long and slender, or so short they resemble cocci)
Spirilla: resemble a corkscrew
What are three other ways to identify bacteria besides morphology?
Presence of flagella (motility)
Formation of endospores
How is motility in bacteria tested?
Bacteria are injected into a tube containing a dye that turns red when oxidized by growing bacteria- distribution of red dye indicates swimming ability
How is formation of endospores determined in bacteria?
Sample can withstand extreme conditions (high temperatures) and will grow at optimal conditions
What are the different enzymatic activities in bacteria that can be tested?
How can resolution be decreased (improved)?
Using an illumination source with a smaller wavelength
Increasing the numerical aperture of the objective lens, as well as using immersion oil with the 100x objective lens
How is the resolution value calculated?
Why do Gram-positive bacteria stain purple?
They contain a thick layer of peptidoglycan which retains the iodine-crystal violet complex CV-I, causing cells to hold the dye and thus retain the purple colour even after the ethanol wash
Why do Gram-negative bacteria stain pink?
Gram negative cells have only a thin layer of peptidoglycan, so the CV-I complex easily washes out, and thus the cells are able to be counterstained pink because they no longer contain any crystal violet dye
What is the effect of high temperatures on bacteria that do not form endospores?
High temperatures will kill cells that do not produce endospores as it can damage cell membranes and denature proteins, resulting in cells that are unable to function
Microtubules, microfilaments, and intermediate filaments which function in cell structure, cell motility (flagella and cilia - microtubules as part of their ultrastructure), and various biological processes
Mitochondria, chloroplasts, Golgi apparatus, endoplasmic reticulum, nucleus, and vacuoles/vesicles
Function in cell motility (flagella and cilia) and one organism also uses cilia to propel food towards its oral groove
Used in amoeboid movement and organelle movement (intermediate filaments give nucleus its shape - nuclear lamina)
How do microfilaments act in amoeboid movement?
Phase contrast microscope
Phase contrast microscope
Cells and medium have different refractive index, and therefore light traveling through a material with different refractive indexes show a change in phase of light waves, which the microscope then translates to a change in light intensity (areas of higher refractive index appear darker)
Some cell structures contain autofluoresce, some require staining - by exposing cells to several stains at once, different structures will fluoresce different colours
How can contrast be improved?
Staining compounds (vital stain - living cells/tissue, or dead cells and tissues), using special types of microscopes to manipulate light, and by reducing the amount of light
What effect would cytochalasin (inhibits microtubules) and cholchicine (inhibits microfilaments) have on Pelomyxa?
Cytochalasin would cause amoeba to become sessile, as the microfilaments are responsible for the movement of amoebas - cholchicine would have no effect
What effect would cytochalasin (inhibits microtubules) and cholchicine (inhibits microfilaments) have on Euglena?
What effect would cytochalasin (inhibits microtubules) and cholchicine (inhibits microfilaments) have on motile prokaryotes?
What is phagocytosis, and how does it differ from receptor-mediated endocytosis?
Phagocytosis is the process that Paramecium use to take in food. It differs from receptor-mediated endocytosis in that receptor-mediated endocytosis is very specific and allows the cell to acquire bulk quantities of specific substances, whereas phagocytosis is more general and can take in different substances
Organelles (enzymes) that digest or break down waste materials and cellular debris, such as worn out organelles, food particles, and engulfed bacteria and viruses
Contains cell's genetic information, and is surrounded by the nuclear lamina (made up of intermediate filaments) to protect the DNA
Engulfs food via phagocytosis and uses lysosomes to break the food down
Organelles found in plant cells that are used in photosynthesis - they capture light and convert it to usable energy
Proteins that catalyze metabolic reactions without being consumed or destroyed by the molecule - lower a reaction's activation energy (substrate specific)
Molecule to be reacted, that fits into a uniquely shaped pocket of the enzyme called the active site and binds with the enzyme as it is converted into the end product
Allows plants to use starch it has stored after photosynthesis - takes amylose and breaks it down into smaller molecules by hydrolysis (glucose molecules, maltose, and shorter chains of amylose)
Polymeric macromolecule composed of glucose monomers that is too large to pass through a cell membrane
What are enzymes made up of?
Enzymes are proteins, which are made up of amino acids
What is an active site?
An active site is a site that uniquely fits the substrate specific to the enzyme, and will activate the enzyme once the substrate binds to the site
Energy transfer from light to chemical bonds through series of light reactions
What series of reactions occurs in photosynthesis?
Light energy from sun strikes pigments in thylakoid membrane of chloroplast, which is transformed into excited electrons (electrical energy), then into chemical energy in the form of bonds in ATP and NADPH molecules, and the ATP and NADH molecules are then used to power the fixation of carbon dioxide into sugar molecules (Calvin cycle, occurring in stroma)
Determines which wavelengths of light the chloroplasts maximally absorb (these wavelengths also produce the highest rates of photosynthetic activity)
Chloroplast suspension is mixed with indicator dye DCPIP - as DCPIP accepts electrons from the electron transport chain of Photosystem II it becomes reduced and therefore colourless, allowing absorptions to be measured to determine concentrations
What were the controls in the photosynthesis in spinach chloroplasts experiment?
Controls must show that DCPIP is stable and there is no other source of electrons to reduce DCPIP (colour does not change spontaneously) and that the colour of chloroplast suspension is stable and does not change colour spontaneously
What are the independent and dependent variables for the absorbance spectrum of photosynthesis in spinach chloroplasts?
The independent variable is the wavelength of light, and the dependent variable is the absorbance of the chloroplast suspension at varying wavelengths of light
What are the independent and dependent variables for the action spectrum of photosynthesis in spinach chloroplasts?
The independent variable is the colour of light, and the dependent variable is the absorbance of the solution
Why are spinach leaves green?
Spinach leaves are green because they maximally absorb blue and red wavelengths, and green wavelengths are absorbed the least and are therefore reflected back the most, resulting in the green colour
At what wavelength is the action spectrum measured at?
The wavelength that photosynthesis occurs at maximally
Where do light reactions take place in the chloroplast? Reactions of the Calvin cycle?
Light reactions take place in the stroma of the chloroplast, and reactions of the Calvin cycle take place in the stroma
Measuring oxygen levels, sugar produced, or carbon dioxide levels
How is ATP produced?
Through the catabolism of carbohydrates, proteins, and fats
Glycolysis in eukaryotes
Cytosolic reactions to convert glucose to pyruvate (one molecule of glucose results in the net production of 2 molecules of ATP via substrate level phosphorylation
After glycolysis, what occurs in the presence of oxygen?
Eukaryotes convert pyruvate into acetyl CoA, which is transported to Kreb's cycle in mitochondria (produces 2 more molecules of ATP) and then oxidative phosphorylation (the transfer of electrons from food to oxygen) produces the rest of the ATP molecules (carbon dioxide is also formed as a by-product)
After glycolysis, what occurs in the absence of oxygen?
Pyruvate is degraded via a series of cytostolic pathways - lactic acid fermentation and alcohol fermentation (produces ethanol and carbon dioxide, regenerates NAD+ - required for glycolytic pathway)
What sort of feedback system occurs in alcohol fermentation?
Fermentation of glucose produces ethanol, but high concentrations of ethanol are toxic to yeast
Physiological response curve
Why is fermentation necessary?
Why was the yeast flask swirled prior to adding yeast to each tube?
To re-suspend the yeast and therefore the ensure that similar concentrations of yeast were present in each tube (constant)
What would happen if the metabolism in yeast experiment were done without the 10 minute pre-incubation period?
The lag phase of the physiological response curve would be significantly longer as the pre-incubation period brings the tube to a temperature at which yeast metabolizes glucose most effectively, and therefore without the incubation period the yeast would not metabolize glucose as well
What process are the yeast in the Durham tube undergoing?
The yeast are undergoing fermentation - other eukaryotes undergo aerobic respiration, and only prokaryotes undergo anaerobic respiration
What metabolic processes occur in the cytoplasm?
Alcohol fermentation, ATP production, glycolysis, and NADH production
What metabolic processes occur in the mitochondria?
ATP production, Krebs cycle, electron transport chain, and NADH production
Bacterial genomic DNA
Consists of a double stranded DNA helix arranged in a circle that is anchored to the bacterial plasma membrane - 4000 genes that encode all the functions of the bacterial cell
Bacterial plasmid DNA
Floats freely in cytoplasm of bacterial cell
Circular and can assume supercoiled conformation in which circular double helix molecule twists on itself
Much smaller than genomic DNA (2- 25 genes)
Can sometimes conform extra properties to the cell that allow the cell to survive in conditions that it could not survive without the plasmid DNA (only when there is selective pressure)
Arranged in linear strands (chromosomes - 23 pairs) in nucleus of cell (30 000 - 35 000 genes - high molecular weight DNA)
Can be used to analyze small amount of plasmid DNA - DNA is not very pure and maxi prep must be used for further analysis as it is a larger quantity of very pure DNA - separates plasmid DNA from bacterial genomic DNA based on size and conformation
How can high molecular weight (HMW) DNA be extracted?
High affinity for glass - buffer solution must contain Tris and EDTA as it binds magnesium ions which are required for DNAse, preventing DNAse from functioning and degrading the DNA into nucleotides
What does centrifuging do?
Creates a centrifugal force that causes bacterial cells to collect in a pellet at the bottom of the tube - liquid above is referred to as the supernatant
What does vortexing do?
Vortexing disrupts the pellet of cells so that they may be re-suspended
Why is STE added to the DNA treatment?
Washes the medium away from the cells
What is Solution I in the DNA treatment?
A buffered, isotonic solution that is used to re-suspend bacterial cells
What is Solution II in the DNA treatment?
Contains sodium dodecylsulfate (SDS) and sodium hydroxide (an alkali) - SDS denatures proteins and disrupts the plasma membrane, causing the cell to lyse and releasing cell components into the solution, and NaOH raises the pH of the lysate to denature the hydrogen bonds between the base pairs of DNA, separating the helix
What is Solution III in the DNA treatment?
Acidic potassium acetate solution that neutralizes the pH in the lysate so that some hydrogen bonds in the DNA will re-form in random base pairs, resulting in a tangled, insoluble mass of DNA - hydrogen bonds in the plasmid DNA reform between the original complementary base pairs (when solution is placed on ice potassium forms white, insoluble mass with SDS that precipitates out along with many of the proteins, cell wall, debris, and genomic DNA)
What does centrifuging do to the genomic DNA-potassium-SDS-protein-cell wall complex?
Causes the complex to pellet in the bottom of the tube and the plasmid to remain in the supernatant solution
What does the 95% ethanol wash do in the DNA treatment?
Removes water molecules from macromolecules by decreasing hydrogen bonding between water molecules and macromolecules (plasmid DNA and RNA come out of solution and precipitate, so that they may be centrifuged into a pellet)
What does the 70% ethanol wash do in the DNA treatment?
Removes the salts which were not removed with the 95% ethanol, and hydrates the pellets slightly so that it may dissolve in the aqueous solution
Why is the 30 minute incubation period necessary in the DNA extraction?
Why is sodium acetate used to precipitate HMW DNA?
The salt ions compete with macromolecules (DNA) for the water molecules
What are the 2 properties of DNA that allow you to separate genomic DNA from plasmid DNA
Size - big genomic DNA precipitates faster with centrifugation
Conformation (shape of molecule) - supercoiled plasmid DNA maintains its shape even when hydrogen bonds in backbone are broken
How can HMW DNA be extracted from solution?
Its affinity for glass and the fact that it forms very long "threads" of DNA
What would happen if the tube were vortexed after the addition of Solution II?
Genomic DNA would break and would not all be centrifuged out, therefore contaminating plasmid DNA
What is the difference between genomic DNA, plasmid DNA, and eukaryotic DNA?
Genomic DNA contains the majority of genes needed for the bacterial cell to function
Plasmid DNA is a small, circular structure of DNA in the cell cytoplasm that contains genes that can allow the bacteria to survive in conditions where it could otherwise not survive
Eukaryotic DNA is much larger and is contained within the nucleus, in 23 pairs of chromosomes, encoding all the genes necessary for the survival of the eukaryote
What is the purpose of the 95% ethanol and 70% ethanol wash?
95% ethanol dehydrates the cell
70% ethanol treatment removes the salts and rehydrates the plasmid DNA, allowing it to dissolve faster
Why must plasmid DNA be kept on ice following incubation?
DNAse will break down DNA at room temperature - T solution has EDTA to inactivate DNAse
What was the experiment performed by Avery, MacLeod, and McCarty?
Tested various cellular macromolecules for their ability to transform non-virulent Streptococcus pneumoniae into virulent bacteria - discovered DNA was the only macromolecule capable of transforming non-virulent bacteria into virulent bacteria
How did the experiment performed in Biol 107 differ from Avery et al?
Escherichia coli was used
E. coli was examined for transformation by a gene on the plasmid DNA instead of the genomic DNA
E. coli cells needed to be made competent to uptake DNA using a calcium chloride solution
Mice were not used (medium containing kanamycin was used instead)
Only DNA was focused on (as opposed to various parts of the cell)
How were kanamycin sensitive E. coli cells made proficient to take up DNA?
A calcium chloride solution was made, which created holes in the cellular membranes (competent cells)
How is a competent cell transformed?
If plasmid DNA entering competent cells is capable of replicating, the competent cells will be genetically altered or transformed (kanamycin resistant) - all descendants of transformed cells should be genetically altered
What is kanamycin?
Antibiotic belonging to the family of antibiotics characterized by their ability to inhibit protein synthesis in prokaryotic cells - they are transported into the cell by oxygen dependent active transport system and irreversibly inhibit protein synthesis by binding to a small subunit of ribosomes in bacterial cell, so cells are unable to synthesize proteins - cell death
What occurs in kanamycin resistant cells?
Phosphotransferase enzyme is encoded and expressed in the presence of kanamycin, which phosphorylates (adds a phosphorous group) to kanamycin and renders the antibiotic inactve
Use of a DNA template to synthesize RNA
Reading of mRNA to produce protein
Plate Count Method
Viable cell count (living cells only) in which original cell suspension is diluted into suspensions of decreasing cell concentration, which are spread onto the surface of an agar medium and allowed to incubate so that single cells may grow into a colony - following incubation colonies may be counted, and each is representative of a single cell originally deposited on the plate
Petroff-Hausser Counting Chamber
Total cell count (living and dead) using a specially designed microscope slide with a depressed surface and etched grid, where a thin layer of cell suspension of known volume is spread and the number of cells in the volume is directly counted with the aid of a microscope
Optical Density (OD)
Indirect method of total cell count, measuring turbidity (cloudiness of a solution due to the presence of particles such as cells), measured using a spectrophotometer, and developing a standard curve
Why is Solution T (Tris) buffer used in the Transformation of Bacterial Cells lab?
Maintains the pH at 8.0 and is the solvent for plasmid DNA
Solution B is the solution used to dissolve DNAse but does not contain DNAse
What occurs during the first incubation period of the Transfomation of Bacterial Cells lab?
What occurs during the heat shock incubation period in the Transformation of Bacterial Cells lab?
Helps the plasmid DNA enter the competent cells and induces the expression of survival genes necessary to repair damage to the plasma membrane
What does the third incubation period in the Transformation of Bacterial Cells lab do?
Allows time for kanamycin resistance gene to be expressed - must be transcribed into mRNA, then mRNA must translate it into a polypeptide chain (phosphotransferase)
The plate which does not contain plasmid DNA, and instead contains solution B, Tris buffer, and competent cells
The plate that contains plasmid DNA that has been broken down into nucleotides by DNAse, as well as solution B and competent cells
There would be colony growth on plate 5+K, as the DNAse cannot enter the competent cells and the plasmid DNA would not have been broken down - the kanamycin resistant gene would have been expressed in the competent cells
What would occur if the environment in which the E. coli was grown was anaerobic?
The kanamycin would not affect the growth of the cells, as kanamycin enters the cell in an oxygen dependent manner
About this deck
Size: 132 flashcards |
“Writing is a way of talking without being interrupted.” — Jules Renard
“ideas to get your student’s pencils moving”
Writing is and isn’t an easy thing to do in the classroom. Especially nowadays when students don’t have long attention spans and are more and more “digital” and visual learners.
However, it is a vital skill that opens up a world of possibilities for any student. Written communication in whatever language, even with the advent of the internet, is still a necessity. Writing allows communication, controlled and deliberate – POWERFUL, communication. So we have to get our students writing more and better. How?
Below, find a rundown of what I consider the “standard” writing activities for any age group. Just change the topic/theme. Most are for any classroom, EFL / ESL or the regular classroom. My belief is that writing in English is writing in English. Whether it be a second language or first makes no difference because the “eating” is all the same.
I’ve divided the activities into different categories. These are just for the sake of having some kind of organization. I’ve also labeled them
WUP – for a warm up writing activity and something to do quickly.
CP – Controlled practice. Writing activities that help the beginning writer and offer support, repetition and guidance.
F – Free writing activities which activate student learning and allow them to practice what they already know and “test the waters” so to speak.
Where appropriate, I’ve linked to some resources that compliment the writing activity as described.
Listen — Write
There are many ways to “spice” up the standard dictation. The simplest is to have the students fold a blank piece of paper “hamburger” style (Up/down) 4 times. Unfold and they have a nice 8 line piece of paper. Speak 8 sentences , repeating each several times as the students write. Get the students to record their answers on the board and correct. Collect and keep in a portfolio!
There are many online sites where students can do the same but in a computer lab or at home. Or the teacher can even try in the classroom.
http://www.listen-and-write.com/audio – for older students
http://www.learner.org/interactives/spelling/ – for young learners
2. Story Rewriting
The teacher reads a story or the class listens to an audio story. After, students make a story board (just fold a blank page so you have 8 squares) and draw pictures. Then, they write the story based on those pictures. Very simple and powerful! – F
The students close their eyes and the teacher describes a scene. Play some nice background music. The students then write and describe the scene they imagined, sharing their scene afterwards with the class or a classmate.
4. Pop Song Rewrite
Play a familiar pop song. One with a “catchy” chorus. Afterwards, write out the chorus on the board with some of the words missing. Students can then rewrite the chorus and sing their own version. Higher level students can simply write their own version without help. Here’s a very simple example –
He’s got the whole world in his hands
He’s got ___________ and ___________
In his hands. (3x)
He’s got the whole world in his hands
Watch —– Write
Students watch a TV Commercial. Then, they write their own script based on that commercial but focused on a different product. Afterwards, they can perform. F
2. Short videos.
Just like a story but this time students watch. Then, they can rewrite / respond / reflect. Students can choose to reflect on one standard Reading Response question or as part of a daily journal. Ex. The best part was ….. / If I had made the video, I would have ……
Short videos are powerful and if well chosen can really get students writing in a reflective manner. CP / F
How to Videos
Students can watch a short “How to” video that describes a process. There are some excellent sites with User Generated Content. Expert Village and eHow are recommended. After the students watch the video several times, they can write out the steps using transitions which the teacher lists on the board. [First, first off, To begin, then, after that, next, most importantly, finally, last but not least, to finish ] CP / F
3. Newscasts / Weather reports
Watch the daily news or weather report. Students write in groups or individually, their own version of the news for that week/day. Then perform for the class like a real news report! F
4. Travel Videos
Watch a few travel videos (there are many nice, short travel “postcard” videos online). Groups of students select a place and write up a report or a poster outlining why others should visit their city/country. Alternately, give students a postcard and have them write to another student in the classroom as if they were in that city/country. For lower leveled students, provide them with a template and they just fill in the details. Ex.
I’m sitting in a ………… drinking a …………… I’ve been in ……. for ………. days now. The weather has been
……………. Yesterday I visited the ………….. and I saw …………….. Today, I’m going to ………………. I highly recommend ……………….. See you when I get home …………………
CP / F
Look —– Write
1. Pictures / Slideshows
Visuals are a powerful way to provide context and background for any writing. Make sure to use attractive, stimulating and if possible “real” photos to prompt student writing. Students can describe a scene or they can describe a series of pictures from a slideshow.
An excellent activity is to show a nice photo and get students to “guess” and write their guess in the form of the 5Ws. They answer all the 5w questions and then share their thoughts with the class.
Show a picture and get students to write a story or use it as background for a writing prompt. For example, Show a picture of a happy lottery winner. Ask students to write in their journal – If I won a million dollars I would ……
This is a much better way to “prompt” writing than simple script! – CP
Show students a selection of fairly similar pictures. The students describe in writing one of the pictures (faces work really well). They read and the other students listen and “guess” which picture is being described. Similar to this listening activity. CP
Provide students with a series of pictures which describe a story. I often use Action Pictures. Students write about each picture, numbering each piece of writing for each picture. The teacher can guide lower level students like this Mr. X’s Amazing Day example. After editing, the students cut up the pictures and make a storybook. Gluing in the pictures, coloring, decorating and adding their own story text. Afterwards read to the whole class or share among the class. CP / F
Provide students with a sequence of pictures which are scrambled. The students must order the pictures and then write out the process. Ex. Making scrambled eggs. F
Read —- Write
1. Reading Journal / Reading Response
The students read a story and then respond by making a reflective journal entry. Alternatively, the students can respond to a reading response question like, “Which character did you like best? Why?” F
Read a short story and then give students a copy of the story with some text missing. The students can fill it in with the correct version OR fill it in and make the story their own.
These are stories where words are replaced with icons/pictures. Students can read the story and then write out the whole story, replacing the pictures with the correct text. Here are some nice examples. – CP
3. Opinion / Essay
Select an article or OP Ed piece that students would find interesting or controversial. After reading and discussing, students can respond with a formal essay or piece of writing reflecting their opinion. Read them anonymously afterwards and get the class to guess who wrote it! F
4. Giving Advice
Students read a problem provided by the teacher (even better, get students to provide the problem by having them write down what they need advice on). This can often be an Ann Landers style request for advice from a newspaper. Students write their own response, giving advice. F
5. Running dictation
This is a lot of fun but quite noisy. Put students into groups of 3 or 4. For each group, post on the wall around the classroom, a piece of writing (maybe a selection of text you will be reading in your lesson). One student is appointed as the secretary. The other students must “run” to where their piece of writing is on the wall and read it. Then run back and dictate it to the secretary who records it. Continue until one group is finished (but check that they got it right!). CP
Think —- Write
1. Graphic Organizers
These you can make on your own by having students draw and fold blank sheets of paper or by giving them a pre-designed one. Students write out their thoughts on a topic using the organizer. An alphabet organizer is also an excellent activity in writing for lower level students. Graphic organizers and mind maps are an excellent way “first step” to a longer writing piece and are an important pre-writing activity. WUP
2. Prompts / Sentence Starters
Students are prompted to finish sentences that are half started. They can write X number of sentences using the sentence starter. Many starters can be found online. Prompts are also an excellent way to get students thinking and writing. Every day, students can “free write” a passage using the daily prompt (ex. What I did this morning etc… ) Creative writing of this sort really motivates students to write. There are many lists online you can use.
3. Thinking Games
Using a worksheet, students play the game while writing down their responses in grammatical sentences. What the Wordle / Not Like the Other and Top 5 are some games I’ve made and which help students begin to write. Each has a worksheet which students fill out. CP
4. Decoding / Translating
Translating a passage into English can be a good writing activity for higher level students.
Students love their cell phones and Transl8it.com is a handy way to get students interested in writing. Simply put in English text and Transl8it.com will output “text messaging”. Give this to students to decode into standard English and then check against the original. Lots of fun! See the games I’ve designed (Pop Song / Dialogues ) using this principle of decoding text messaging. CP
5. Forms / Applications
Students need to practice writing that will be of use to them directly in the wider world. Forms and filling in applications are a valuable way to do this. Fill in one together as a class and then get students to do this same for themselves individually. – CP
6. Journals / Reflection / Diaries
This type of free writing activity should be done on a regular basis if used in class. Use a timer and for X minutes, students can write upon a topic that is important to them, that day. Alternatively, students can write at the end of the day and record their thoughts about the lesson or their own learning. These are all excellent ways for the teacher to get to know their students. One caution – don’t correct student writing here! Comment positively on the student’s writing – the goal is to get them feeling good about writing and “into” it. – F
7. Tag Stories / Writing
Students love this creative exercise. Fold a blank piece of paper vertically (Hamburger style) 4 times. You’ll have 8 lines. On the first line, students all write the same sentence starter. Ex. A man walked into a bank and ……..
Next, students finish the sentence and then pass their paper to the student on their left/right. That student reads the sentence and continues the story on the next line. Continue until all 8 lines are completed. Read the stories as a class – many will be hilarious! I often do this with a “gossip” variation. I write some gossip “chunks” on the board like; “I heard that..” , “I was told…” “The word on the street is…” “Don’t pass it around but…”. Students choose one and write some juicy gossip about the student to their right. They then pass their paper to the left with everyone adding onto the gossip. Students really get into this! CP / F
8. Describe and guess
Students think of a person / a place or a thing. They write a description of them / it and they are read out and others students guess.
Jokes and riddles are also effective for this. Students write out a joke or riddle they know and then they are read and other students try to guess the punchline. – F
TEXT —– Write
1. Sentence Chains
The teacher writes a word on the board and then students shout out words that follow using the last letter(s). The more last letters they use, the more points they get. The teacher keeps writing as quick as possible as the students offer up more correct words. Ex. Smilengthosentencementality…..
Give students a blank piece of paper and in pairs with one student being the secretary, they play! This is a great game for simple spelling practice and also to get students noticing language and how words end/begin. They can also play for points. Compound words and phrases are acceptable! – WUP
2. Guided Writing
This is a mainstay of the writing teacher’s toolkit. Students are either given a “bank” of words or can write/guess on their own. They fill in the missing words of a text to complete the text. Take up together and let students read their variations. A nice adaptation to guided writing for lower level students is for them to personalize the writing by getting them to draw a picture for the writing passage to illustrate and fortify the meaning. Here’s a nice example. CP
Use a time line to describe any event. Brainstorm as a class. Then students use the key words written on the board, to write out the time line as a narrative. Really effective and you can teach history like this too! Biographies of individuals or even the students themselves are a powerful writing activity and timelines are a great way to get them started. – F
Students are given notes (the classic example is a shopping list but it might be a list of zoo animals / household items etc…) and then asked to write something using all the noted words. This usually focuses on sequence (transitions) or location (prepositions). F
5. Grammar Poems
Grammar poems are short poems about a topic that students complete using various grammar prompts. This form of guided writing is very effective and helps students notice various syntactical elements of the language.
Put the grammar poem on the board with blanks. Here are some examples but it could be on any topic (country, famous person, my home, this school, etc..). Fill out as a class with one student filling it in. Then, students copy the poem and complete with their own ideas. Change as needed to stress different grammatical elements. And of course, afterwards SHARE. Present some to the class and display on a bulletin board. Your students will be proud of them!
SPEAK — Write
1. Surveys / Reports
Students have a survey question or a questionnaire. They walk around the class recording information. After, instead of reporting to the class orally, they can write up the report about their findings.
This can also be used with FSW (Find Someone Who) games. Students use a picture bingo card to walk around the classroom and ask students yes/no questions. They write the answers with a check or X and the student’s name in the box with the picture. After, they write up a report about which student ……. / didn’t …… certain things. CP
2. Reported Speech
Do any speaking activity or set of conversation questions. Afterwards, students report back by writing using reported speech, “ Susan told me that she ………..” and “ Brad said that ………..” etc….. CP
3. Introducing each other
Students can interview another classmate using a series of questions / key words given by the teacher. After the interview of each other is over, students can write out a biography of their partner and others can read them in a class booklet. – F
4. In class letter writing
Writing for a purpose is so important and nothing makes this happen better than in class letter writing. Appoint a postman and have each student make a post office box (it could just be a small bag hanging from their desk). The students can write each other (best to assign certain students first) and then respond to their letter. Once it gets started, it just keeps going and going… – F
4. Email / messaging / chat / social networking
This is an excellent way to get students speaking by writing. Set up a social networking system or a messaging / emailing system for the students. They can communicate and chat there using an “English only” policy. Use videos / pictures like in class – to promote student discussion and communication. Projects online foster this kind of written communication and using an CMS (Content Management System) like moodle or atutor or ning can really help students write more. – F
5. Class / School English newspaper or magazine
Students can gain valuable skills by meeting and designing a school English newsletter. Give each student a role (photographer, gossip / news / sports / editor in chief / copy editor etc…) and see what they can do. You’ll be surprised! – F
WRITE —— Do
Students can write dialogues for many every day situations and then act them out for the class. The teacher can model the language on the board and then erase words so students can complete by themselves and in their own words. Here’s a neat example using a commercial as a dialogue. – CP
Students draw a picture and then write a description of the picture. They hand their description to another student who must read it and then draw the picture as they see it. Finally, both students compare pictures! – F
3. Tableaus / Drama
Students write texts of any sort. Then the texts are read and other students must make a tableau of the description or act out the text in some manner. For example – students can write about their weekend. After writing, the student reads their text and other students act it out or perform a tableau. F
4. Don’t speak / Write!
I once experimented with a class that wouldn’t speak much by putting a gag on myself and only writing out my instructions. It worked and this technique could be used in a writing class. Students can’t speak and are “gagged”. Give them post it notes by which to communicate with others. Instruct using the board. There are many creative ways to use this technique! – F
RECOMMENDED BOOKS 4 TEACHERS
I highly recommend the following two books for ideas and some general theory on how to teach writing. Purchase them for reference.
1. HOW TO TEACH WRITING – Jeremy Harmer
Very insightful and cleanly, simply written. The author explores through example and description, all the facets and theory behind that “looking glass” which we call teaching. I use this as a course text for my methodology class for in-service teachers.
2. Oxford Basics: Simple Writing Activites
- Jill and Charles Hadfield
This book (and series) is a gem! Jill Hadfield knows what working EFL / ESL teachers need and in this book there are 30 simple writing activities which teachers can use with a wide variety of levels and with only a chalkboard and a piece of chalk / paper.
See my Blog post and download the list of my TOP 10 WRITING WEBSITES FOR TEACHERS AND STUDENTS |
A long time ago I wrote the article The Dull Case of Emissivity and Average Temperatures and expected that would be the end of the interest in emissivity. But it is a gift that keeps on giving, with various people concerned that no one has really been interested in measuring surface emissivity properly.
All solid and liquid surfaces emit thermal radiation according to the Stefan-Boltzmann formula:
E = εσT4
where ε=emissivity, a material property; σ = 5.67×10-8 ; T = temperature in Kelvin (absolute temperature)
and E is the flux in W/m²
More about this formula and background on the material properties in Planck, Stefan-Boltzmann, Kirchhoff and LTE.
The parameter called emissivity is the focus of this article. It is of special interest because to calculate the radiation from the earth’s surface we need to know only temperature and emissivity.
Emissivity is a value between 0 and 1. And is also depends on the wavelength of radiation (and in some surfaces like metals, also the direction). Because the wavelengths of radiation depend on temperature, emissivity also depends on temperature.
When emissivity = 1, the body is called a “blackbody”. It’s just the theoretical maximum that can be radiated. Some surfaces are very close to a blackbody and others are a long way off.
Note: I have seen many articles by keen budding science writers who have some strange ideas about “blackbodies”. The only difference between a blackbody and a non-blackbody is that the emissivity of a blackbody = 1, and the emissivity of a non-blackbody is less than 1. That’s it. Nothing else.
The wavelength dependence of emissivity is very important. If we take snow for example, it is highly reflective to solar (shortwave) radiation with as much as 80% of solar radiation being reflected. Solar radiation is centered around a wavelength of 0.5μm.
Yet snow is highly absorbing to terrestrial (longwave) radiation, which is centered around a wavelength of 10μm. The absorptivity and emissivity around freezing point is 0.99 – meaning that only 1% of incident longwave radiation would be reflected.
Let’s take a look at the Planck curve – the blackbody radiation curve – for surfaces at a few slightly different temperatures:
The emissivity (as a function of wavelength) simply modifies these curves.
Suppose, for example, that the emissivity of a surface was 0.99 across this entire wavelength range. In that case, a surface at 30°C would radiate like the light blue curve but at 99% of the values shown. If the emissivity varies across the wavelength range then you simply multiply the emissivity by the intensity at each wavelength to get the expected radiation.
Sometimes emissivity is quoted as an average for a given temperature – this takes into account the shape of the Planck curve shown in the graphs above.
Often, when emissivity is quoted as an overall value, the total flux has been measured for a given temperature and the emissivity is simply:
ε = actual radiation measured / blackbody theoretical radiation at that temperature
[Fixed, thanks to DeWitt Payne for pointing out the mistake]
In practice the calculation is slightly more involved, see note 1.
It turns out that the emissivity of water and of the ocean surface is an involved subject.
And because of the importance of calculating the sea surface temperature from satellite measurements, the emissivity of the ocean in the “atmospheric window” (8-14 μm) has been the subject of many 100′s of papers (perhaps 1000′s). These somewhat overwhelm the papers on the less important subject of “general ocean emissivity”.
Aside from climate, water itself is an obvious subject of study for spectroscopy.
For example, 29 years ago Miriam Sidran writing Broadband reflectance and emissivity of specular and rough water surfaces, begins:
The optical constants of water have been extensively studied because of their importance in science and technology. Applications include a) remote sensing of natural water surfaces, b) radiant energy transfer by atmospheric water droplets, and c) optical properties of diverse materials containing water, such as soils, leaves and aqueous solutions.
In this study, values of the complex index of refraction from six recent articles were averaged by visual inspection of the graphs, and the most representative values in the wavelength range of 0.200 μm to 5 cm were determined. These were used to find the directional polarized reflectance and emissivity of a specular surface and the Brewster or pseudo-Brewster angle as functions of wavelength.
The directional polarized reflectance and emissivity of wind-generated water waves were studied using the facet slope distribution function for a rough sea due to Cox and Munk .
Applications to remote sensing of sea surface temperature and wave state are discussed, including effects of salinity.
Emphasis added. She also comments in her paper:
For any wavelength, the total emissivity, ε, is constant for all θ [angles] < 45° [from vertical]; this follows from Fig. 8 and Eq. (6a). It is important in remote sensing of thermal radiation from space, as discussed later..
The polarized emissivities are independent of surface roughness for θ < 25°, while for θ > 25°, the thermal radiation is partly depolarized by the roughness.
This means that when you look at the emission radiation from directly above (and close to directly above) the sea surface roughness doesn’t have an effect.
I thought some other comments might also be interesting:
The 8-14-μm spectral band is chosen for discussion here because (a) it is used in remote sensing and (b) the atmospheric transmittance, τ, in this band is a fairly well-known function of atmospheric moisture content. Water vapor is the chief radiation absorber in this band.
In Eqs. (2)-(4), n and k (and therefore A and B) are functions of salinity. However, the emissivity value, ε, computed for pure water differs from that of seawater by <0.5%.
When used in Eqs. (10), it causes an error of <0.20°C in retrieved Ts [surface temperature]. Since ε in this band lies between 0.96 and 0.995, approximation ε= 1 is routinely used in sea surface temperature retrieval. However, this has been shown to cause an error of -0.5 to -1.0°C for very dry atmospheres. For very moist atmospheres, the error is only ≈0.2°C.
One of the important graphs from her paper:
Click to view a larger image
Emissivity = 1 – Reflectance. The graph shows Reflectance vs Wavelength vs Angle of measurement.
I took the graph (coarse as it is) and extracted the emissivity vs wavelength function (using numerical techniques). I then calculated the blackbody radiation for a 15°C surface and the radiation from a water surface using the emissivity from the graph above for the same 15°C surface. Both were calculated from 1 μm to 100 μm:
The “unofficial” result, calculating the average emissivity from the ratio: ε = 0.96.
This result is valid for 0-30°C. But I suspect the actual value will be modified slightly by the solid angle calculations. That is, the total flux from the surface (the Stefan-Boltzmann equation) is the spectral intensity integrated over all wavelengths, and integrated over all solid angles. So the reduced emissivity closer to the horizon will affect this measurement.
Niclòs et al – 2005
One of the most interesting recent papers is In situ angular measurements of thermal infrared sea surface emissivity—validation of models, Niclòs et al (2005). Here is the abstract:
In this paper, sea surface emissivity (SSE) measurements obtained from thermal infrared radiance data are presented. These measurements were carried out from a fixed oilrig under open sea conditions in the Mediterranean Sea during the WInd and Salinity Experiment 2000 (WISE 2000).
The SSE retrieval methodology uses quasi-simultaneous measurements of the radiance coming from the sea surface and the downwelling sky radiance, in addition to the sea surface temperature (SST). The radiometric data were acquired by a CIMEL ELECTRONIQUE CE 312 radiometer, with four channels placed in the 8–14 μm region. The sea temperature was measured with high-precision thermal probes located on oceanographic buoys, which is not exactly equal to the required SST. A study of the skin effect during the radiometric measurements used in this work showed that a constant bulk–skin temperature difference of 0.05±0.06 K was present for wind speeds larger than 5 m/s. Our study is limited to these conditions.
Thus, SST used as a reference for SSE retrieval was obtained as the temperature measured by the contact thermometers placed on the buoys at 20-cm depth minus this bulk–skin temperature difference.
SSE was obtained under several observation angles and surface wind speed conditions, allowing us to study both the angular and the sea surface roughness dependence. Our results were compared with SSE models..
The introduction explains why specifically they are studying the dependence of emissivity on the angle of measurement – for reasons of accurate calculation of sea surface temperature:
The requirement of a maximum uncertainty of ±0.3 K in sea surface temperature (SST) as input to climate models and the use of high observation angles in the current space missions, such as the 55° for the forward view of the Advanced Along Track Scanning Radiometer (AATSR) (Llewellyn-Jones et al., 2001) on board ENVISAT, need a precise and reliable determination of sea surface emissivity (SSE) in the thermal infrared region (TIR), as well as analyses of its angular and spectral dependences.
The emission of a rough sea surface has been studied over the last years due to the importance of the SSE for accurate SST retrieval. A reference work for many subsequent studies has been the paper written by Cox and Munk (1954)..
The experimental setup:
From Niclos (2004)
The results (compared with one important model from Masuda et al 1988):
From Niclos (2004)
Click on the image for a larger graphic
This paper also goes on to compare the results with the model of Wu & Smith (1997) and indicates the Wu & Smith’s model is a little better.
The tabulated results, note that you can avoid the “eye chart effect” by clicking on the table:
Click on the image for a larger view
Note that the emissivities are in the 8-14μm range.
You can see that the emissivity when measured from close to vertical is 0.98 – 0.99 at two different wind speeds.
Konda et al – 1994
A slightly older paper which is not concerned with angular dependence of sea surface emissivity is by Konda, Imasato, Nishi and Toda (1994).
They comment on a few older papers:
Buettner and Kern (1965) estimated the sea surface emissivity to be 0.993 from an experiment using an emissivity box, but they disregarded the temperature difference across the cool skin.
Saunders (1967b, 1968) observed the plane sea surface irradiance from an airplane and determined the reflectance. By determining the reflectance as the ratio of the differences in energy between the clear and the cloudy sky at different places, he calculated the emissivity to be 0.986. The process of separating the reflection from the surface irradiance, however, is not precise.
Mikhaylov and Zolotarev (1970) calculated the emissivity from the optical constant of the water and found the average in the infrared region was 0.9875.
The observation of Davies et al. (1971) was performed on Lake Ontario with a wave height less than 25 cm. They measured the surface emission isolated from sky radiation by an aluminum cone, and estimated the emissivity to be 0.972. The aluminum was assumed to act as a mirror in infrared region. In fact,aluminum does not work as a perfect mirror.
Masuda et al. (1988) computed the surface emissivity as a function of the zenith angle of observed radiation and wind speed. They computed the emissivity from the reflectance of a model sea surface consisting of many facets, and changed their slopes according to Gaussian distribution with respect to surface wind. The computed emissivity in 11 μm was 0.992 under no wind.
Each of these studies in trying to determine the value of emissivity, failed to distinguish surface emission from reflection and to evaluate the temperature difference across the cool skin. The summary of these studies are tabulated in Table 1.
The table summarizing some earlier work:
Konda and his co-workers took measurements over a one year period from a tower in Tanabe Bay, Japan.
They calculated from their results that the ocean emissivity was 0.984±0.004.
One of the challenges for Konda’s research and for Niclòs is the issue of sea surface temperature measurement itself. Here is a temperature profile which was shown in the comments of Does Back Radiation “Heat” the Ocean? – Part Three:
Kawai & Wada (2007)
The point is the actual surface from which the radiation is emitted will usually be at a slightly different temperature from the bulk temperature (note the logarithmic scale of depth). This is the “cool skin” effect. This surface temperature effect is also moderated by winds and is very difficult to measure accurately in field conditions.
Smith et al – 1996
Another excellent paper which measured the emissivity of the ocean is by Smith et al (1996):
An important objective in satellite remote sensing is the global determination of sea surface temperature (SST). For such measurements to be useful to global climate research, an accuracy of ±0.3K or better over a length of 100km and a timescale of days to weeks must be attained. This criterion is determined by the size of the SST anomalies (≈1K) that can cause significant disturbance to the global atmospheric circulation patterns and the anticipated size of SST perturbations resulting from global climate change. This level of uncertainty is close to the theoretical limits of the atmospheric corrections..
It is also a challenge to demonstrate that such accuracies are being achieved, and conventional approaches, which compare the SST derived from drifting or moored buoys, generally produce results with a scatter of ±0.5 to 0.7K. This scatter cannot be explained solely by uncertainties in the buoy thermometers or the noise equivalent temperature difference of the AVHRR, as these are both on the order of 0.2K or less but are likely to be surface emissivity/reflectivity uncertainties, residual atmospheric effects, or result from the methods of comparison
Note that the primary focus of this research was to have accurate SST measurements from satellites.
From Smith et al (1996)
The experimental work on the research vessel Pelican included a high spectral resolution Atmospheric Emitted Radiance Interferometer (AERI) which was configured to make spectral observations of the sea surface radiance at several view angles. Any measurement from the surface of course, is the sum of the emitted radiance from the surface as well as the reflected sky radiance.
- ocean salinity
- intake water temperature
- surface air temperature
- wind velocity
- SST within the top 15cm of depth
There was also independent measurement of the radiative temperature of the sea surface at 10μm with a Heimann broadband radiation thermometer “window” radiometer. And radiosondes were launched from the ship roughly every 3 hours.
Additionally, various other instruments took measurements from a flight altitude of 20km. Satellite readings were also compared.
The AERI measured the spectral distribution of radiance from 3.3μm to 20μm at 4 angles. Upwards at 11.5° from zenith, and downwards at 36.5°, 56.5° and 73.5°.
There’s a lot of interesting discussion of the calculations in their paper. Remember that the primary aim is to enable satellite measurements to have the most accurate measurements of SST and satellites can only really “see” the surface through the “atmospheric window” from 8-12μm.
Here are the wavelength dependent emissivity results shown for the 3 viewing angles. You can see that at the lowest viewing angle of 36.5° the emissivity is 0.98 – 0.99 in the 8-12μm range.
From Smith et al (1996)
Note that the wind speed doesn’t have any effect on emissivity at the more direct angle, but as the viewing angle moves to 73.5° the emissivity has dropped and high wind speeds change the emissivity considerably.
Henderson et al – 2003
Henderson et al (2003) is one of the many papers which consider the theoretical basis of how viewing angles change the emissivity and derive a model.
Just as an introduction, here is the theoretical variation in emissivity with measurement angle, versus “refractive index” as computed by the Fresnel equations:
The legend is refractive index from 1.20 to 1.35. Water, at visible wavelengths, has a refractive index of 1.33. This shows how the emissivity reduces once the viewing angle increases above 50° from the vertical.
The essence of the problem of sea surface roughness for large viewing angles is shown in the diagram below, where multiple reflections take place:
Henderson and his co-workers compare their results with the measured results of Smith et al (1996) and also comment that at zenith viewing angles the emissivity does not depend on the wind speed, but at larger angles from vertical it does.
A quick summary of their model:
We have developed a Monte Carlo ray-tracing model to compute the emissivity of computer-rendered, wind-roughened sea surfaces. The use of a ray-tracing method allows us to include both the reflected emission and shadowing and, furthermore, permits us to examine more closely how these processes control the radiative properties of the surface. The intensity of the radiation along a given ray path is quantified using Stokes vectors, and thus, polarization is explicitly included in the calculations as well.
Their model results compare well with the experimental results. Note that the approach of generating a mathematical model to calculate how emissivity changes with wind speed and, therefore, wave shape is not at all new.
Water retains its inherent properties of emissivity regardless of how it is moving or what shape it is. The theoretical challenge is handling the multiple reflections, absorptions, re-emissions that take place when the radiance from the water is measured at some angle from the vertical.
The best up to date measurements of ocean emissivity in the 8-14 μm range are 0.98 – 0.99. The 8-14 μm range is well-known because of the intense focus on sea surface temperature measurements from satellite.
From quite ancient data, the average emissivity of water across a very wide broadband range (1-100 μm) is 0.96 for water temperatures from 0-30°C.
The values from the ocean when measured close to the vertical are independent of wind speed and sea surface roughness. As the angle of measurement moves from the vertical around to the horizon the measured emissivity drops and the wind speed affects the measurement significantly.
These values have been extensively researched because the calculation of sea surface temperature from satellite measurements in the 8-14μm “atmospheric window” relies on the accurate knowledge of emissivity and any factors which affect it.
For climate models – I haven’t checked what values they use. I assume they use the best experimental values from the field. That’s an assumption. I’ve already read enough on ocean emissivity.
For energy balance models, like the Trenberth and Kiehl update, an emissivity of 1 doesn’t really affect their calculations. The reason, stated simply, is that the upwards surface radiation and the downward atmospheric radiation are quite close in magnitude. For example, the globally annually averaged values of both are 396 W/m² (upward surface) vs 340 W/m² (downward atmospheric).
Suppose the emissivity drops from 0.98 to 0.97 – what is the effect on upwards radiation through the atmosphere?
The upwards radiation has dropped by 4W/m², but the reflected atmospheric radiation has increased by 3.4W/m². The net upwards radiation through the atmosphere has reduced by only 0.6 W/m².
One of our commenters asked what value the IPCC uses. The answer is they don’t use a value at all because they summarize research from papers in the field.
Whether they do it well or badly is a subject of much controversy, but what is most important to understand is that the IPCC does not write papers, or perform GCM model runs, or do experiments – and that is why you see almost no equations in their many 1000′s of pages of discussion on climate science.
For those who don’t believe the “greenhouse” effect exists, take a look at Understanding Atmospheric Radiation and the “Greenhouse” Effect – Part One in the light of all the measured results for ocean emissivity.
On Another Note
It’s common to find claims on various blogs and in comments on blogs that climate science doesn’t do much actual research.
I haven’t found that to be true. I have found the opposite.
Whenever I have gone digging for a particular subject, whether it is the diurnal temperature variation in the sea surface, diapycnal & isopycnal eddy diffusivity, ocean emissivity, or the possible direction and magnitude of water vapor feedback, I have found a huge swathe of original research, of research building on other research, of research challenging other research, and detailed accounts of experimental methods, results and comparison with theory and models.
Just as an example, in the case of emissivity of sea surface, at the end of the article you can see the first 30 or so results pulled up from one journal – Remote Sensing of the Environment for the search phrase “emissivity sea surface”. The journal search engine found 348 articles (of course, not every one of them is actually about ocean emissivity measurements).
Perhaps it might turn out to be the best journal for this subject, but it’s still just one journal.
Broadband reflectance and emissivity of specular and rough water surfaces, Sidran, Applied Optics (1981)
In situ angular measurements of thermal infrared sea surface emissivity—validation of models, Niclòs, Valor, Caselles, Coll & Sànchez, Remote Sensing of Environment (2005)
Measurement of the Sea Surface Emissivity, Konda, Imasato, Nishi and Toda, Journal of Oceanography (1994)
Observations of the Infrared Radiative Properties of the Ocean—Implications for the Measurement of Sea Surface Temperature via Satellite Remote Sensing, Smith, Knuteson, Revercomb, Feltz, Nalli, Howell, Menzel, Brown, Brown, Minnett & McKeown, Bulletin of the American Meteorological Society (1996)
The polarized emissivity of a wind-roughened sea surface: A Monte Carlo model, Henderson, Theiler & Villeneuve, Remote Sensing of Environment (2003)
Note 1: The upward radiation from the surface is the sum of three contributions: (i) direct emission of the sea surface, which is attenuated by the absorption of the atmospheric layer between the sea surface and the instrument; (ii) reflection of the downwelling sky radiance on the sea, attenuated by the atmosphere; and (iii) the upwelling atmospheric radiance emitted in the observing direction.
So the measured radiance can be expressed as:
where the three terms on the right are each of the three contributions noted in the same order.
Note 2: 1/10th of the search results returned from one journal for the search term “emissivity sea surface”:
Remote Sensing of Environment - search results
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The government's human rights record remained poor. Although there were no reports that the government or its agents committed politically motivated killings, security forces committed extrajudicial killings and acted with impunity. There was little political will to address the failure by government authorities to adhere to the rule of law. Detainees were abused, often to extract confessions, and prison conditions were harsh. Human rights monitors reported arbitrary arrests and prolonged pretrial detention, underscoring a weak judiciary and denial of the right to a fair trial. Land disputes and forced evictions, often accompanied by violence, were a continuing problem. The government restricted freedom of speech and the press through the use of defamation and disinformation suits, controlled or influenced the content of television and radio broadcasts, and at times interfered with freedom of assembly. Corruption was endemic and extended throughout all segments of society, including the executive, legislative, and judicial branches of government. Domestic violence and child abuse occurred, education of children was inadequate, and trafficking in women and children persisted. The government offered little assistance to persons with disabilities. Antiunion activity by employers and weak enforcement of labor laws continued, and child labor remained a problem.
In a positive turn, on June 12, the Extraordinary Chambers of the Courts in Cambodia for the Prosecution of Crimes Committed during the Period of Democratic Kampuchea adopted its internal rules to begin prosecuting senior leaders of the Khmer Rouge regime and those most responsible for committing serious crimes. On July 31, the ECCC charged Kaing Guek Eav, alias Duch, with crimes against humanity and subsequently charged four other senior officials; at year's end all were in detention awaiting trial. In addition, on December 10, the government permitted a Human Rights Day march of 500 human rights activists, monks, and other persons and rally of an estimated 2,500 persons in Phnom Penh.
RESPECT FOR HUMAN RIGHTS
Section 1 Respect for the Integrity of the Person, Including Freedom From:
a. Arbitrary or Unlawful Deprivation of Life
There were no reports that the government or its agents committed politically motivated killings. However, human rights nongovernmental organizations (NGOs) reported that extrajudicial killings continued to occur.
The Cambodian Human Rights and Development Association (ADHOC) recorded 53 cases of extrajudicial killings, 14 of which were committed by police, nine by soldiers, six by fishery officials, and the remaining 24 by unidentified government forces. Police arrested perpetrators in four cases.
Political activists continued to be the victims of killings. On February 27, Eang Sok Thoeurn, a Khmer Kampuchea Krom monk, was found dead with his throat cut in the Tronum Chhroeung Monastery in Kandal Province. The deceased monk was discovered the morning after he participated in a demonstration in front of the Vietnamese embassy in Phnom Penh for the rights of Khmer Kampuchea Krom persons living in Vietnam. Police quickly declared the death a suicide and disposed of the body without further investigation. NGOs and Khmer Kampuchea Krom groups suspected the killing was politically motivated.
Active members of political parties were killed during the year, but NGOs and police could not confirm their deaths were politically motivated. On February 14, three unidentified persons killed Sam Rainsy Party (SRP) activist Chea Sovin, spouse of an SRP candidate for the April commune council elections in Battambang Province. On July 27, three unidentified persons shot and killed Kleb Un, SRP commune‑level vice party chairperson in Banteay Meanchey Province. A local police chief reported that the perpetrators fled without robbing the victim or taking anything from the scene. Police arrested one suspect in the case but released him after questioning. In both killings, police stated that investigations continued.
On April 4, police officer Siv Soeun allegedly shot and killed a person he claimed was illegally fishing on private property in the Kompong Siem District of Kampong Cham Province. The victim's family filed a complaint against the police officer but later withdrew the complaint after Siv Soeun allegedly paid the family $3,000 (12 million riel) in compensation. At year's end Siv Soeun had not been charged or arrested.
On November 15, during the eviction of squatters from state land in Choam Ksan commune, Preah Vihear Province, unidentified government forces killed two villagers who protested the eviction. Approximately 150 police, military police, and soldiers evicted 317 families. There was no official investigation into the killings. Police arrested 18 of the squatters, including a deputy governor, on charges of encroachment on state land. The 18 villagers were imprisoned and awaited trial at year's end.
In June the Prey Veng Provincial Court sentenced one suspect in the November 2006 killing of SRP activist Man Meth to 10 years in prison and two others to six months in prison for conspiring in the killing.
On July 18, the Phnom Penh Municipal Court sentenced Heng Pov, former Phnom Penh police commissioner and under secretary of state of the Ministry of Interior (MOI), to an additional 22.5 years in prison for the 2005 illegal detention of a person, use of illegal weapons, and possession of counterfeit currency. Heng Pov was already serving an 18‑year sentence for the 2003 murder of Judge Sok Sethamony, multiple counts of premeditated killings, and involvement in illegal arrests and detentions. During his July trial, Heng Pov stated that Born Samnang and Sok Sam Oeun, the two suspects he ordered arrested in 2005 and who later were convicted for the killing of union activist Chea Vichea, were innocent of the crime. On April 12, the appeals court had upheld 20‑year sentences each for Born Samnang and Sok Sam Oeun. Their lawyers submitted grievances to the Supreme Court, and at year's end they awaited Supreme Court action.
There were no developments in the 2006 cases of SRP activists Koent Chhuon and Thoeung Thear, killed in Preah Vihear and Kampong Cham provinces, respectively. Likewise, there were no developments in the cases of Pao Rum and Khat Thoeun, who died in police custody in Kandal Province in 2006, or the 2006 cases of attempted prison breaks in Kampong Thom and Battambang that left 10 inmates dead. In the case of Nong Sam, who reportedly died June 2006 in a Siem Reap hospital from head injuries received during a beating by police officers, a provincial court prosecutor closed the case, declaring Nong Sam's death a suicide.
There were no developments in the 2005 killings of five SRP activists or in the 2005 case of an attempted escape from Trapoeung Phlong Prison in which 19 prisoners and the prison director were killed. The appeals court took no action in the 2005 deaths of five villagers and injuries to others by government security forces during a mass eviction from disputed land in the village of Kbal Spean in Banteay Meanchey Province.
On February 8, in Prey Veng Province, district‑ and commune‑level deputy police chiefs Bun Samphea, Suos Bunthat, and Hay Chivon, charged in a 2000 killing, failed to appear for their provincial court trial, reportedly stating they were too busy. The court rescheduled the trial to June but postponed it again after the officers said they were too busy to come to the June trial. A new trial date was not set.
Mines dating from the Indochina conflict and Khmer Rouge period continued to cause casualties. According to the Cambodia Mine/UXO Victim Information System, during the year mines and unexploded ordnance caused 63 deaths, 56 amputations, and 222 other injuries.
Vigilante justice and mob killings persisted. ADHOC reported that mobs killed five persons during the year. Few suspects were arrested. In some instances authorities could not protect suspects from angry mobs. NGOs noted that a majority of mob killings were related to thefts, robberies, or suspected witchcraft. On June 9, Yos Chor villagers in Kampong Speu Province killed a person for stealing a neighbor's chickens. On June 21, a mob killed a traditional healer in the Boribo District of Kampong Chhnang Province because they suspected him of witchcraft. Police made no arrests in either case.
On June 27, the Kratie Provincial Court sentenced six persons to sentences ranging from seven to 10 months in prison for the 2006 vigilante justice killing of Sam Roeun. The court convicted them on charges of causing injury, reduced from investigation findings of murder. There were no developments in the February 2006 case of a person beaten and killed for allegedly practicing witchcraft.
On June 30, Khmer Kampuchea Krom monk Tim Sakhorn, head of a pagoda in the Kirivong District of Takeo Province for more than 10 years, disappeared. Previously, on orders of the country's top Buddhist leader, Great Supreme Patriarch Tep Vong, monks from Phnom Penh had defrocked Tim Sakhorn, after which unidentified persons believed to be attached to the MOI pushed him into a vehicle and drove away. The defrocking order stated Tim Sakhorn "broke the solidarity" between Cambodia and Vietnam by using pagodas to spread propaganda that affects the dignity of Buddhism. The monk was known locally for providing food and shelter to Khmer Kampuchea Krom coming from Vietnam. The MOI stated that Tim Sakhorn volunteered to go to Vietnam after he was defrocked, and ministry officials produced a document stating this intent. While signed by Tim Sakhorn, the handwritten document appeared not to be in his writing. On August 2, Tim Sakhorn reappeared in court custody in Vietnam, held on charges of destroying political solidarity. In September the Information Ministry stated that the Cambodian consulate in Ho Chi Minh City was investigating Tim Sakhorn's condition in detention. On November 8, a Vietnamese newspaper reported that a court in Vietnam convicted Tim Sakhorn of undermining solidarity between Cambodia and Vietnam and sentenced him to one year in prison.
On August 10, Land Border Protection Unit 504 soldier Im Bun Ny disappeared in Pailin. According to witnesses, that night four soldiers from his unit invited him to a rubber plantation owned by their unit commander, Brigadier General Pol Sinuon. After Im Bun Ny arrived, the four soldiers beat him and accused him of stealing a gun. Unconfirmed witness reports said Im Bun Ny died from the beating and the soldiers buried his body. At year's end Im Bun Ny was still missing. According to a human rights NGO, local police completed an investigation and submitted findings to the court. The court took no action, and the four soldiers remained at large.
c. Torture and Other Cruel, Inhuman, or Degrading Treatment or Punishment
The constitution such practices; however, beatings and other forms of physical mistreatment of police detainees and prison inmates continued to be a serious problem.
There were credible reports that military and civilian police officials used physical and psychological torture and severely beat criminal detainees, particularly during interrogation. Based on interviews with 1,293 detainees from 18 of the country's 26 prisons, the Cambodian League for the Prevention and Defense of Human Rights (LICADHO) reported that during the year authorities tortured 155 prisoners, of whom 125 were tortured in police custody and 30 in prisons. Kicking, punching, and pistol whipping were the most common methods of physical abuse, but techniques also included electric shocks, suffocation, caning, and whipping with wire. NGOs reported that it was not uncommon for police to torture detained suspects until they confessed to a crime. Courts used forced confessions as legal evidence during trial despite admissibility prohibitions under the law.
NGOs noted that during the year there were 180 cases of physical assaults by local authorities, government agents, or private bodyguards, compared with 164 cases in 2006 and 154 cases in 2005.
On May 27, military police officer Prak Vutha of Phnom Sruoch District, Kampong Speu Province, reportedly arrested Sok Soeun after a small scuffle at a restaurant, kept him in military detention overnight without a warrant, and beat him unconscious. According to ADHOC, Sok Soeun's family gave Prak Vutha two cases of beer in return for Sok Soeun's release. Sok Soeun later filed a complaint with local police that the police did not accept. There was no investigation into the case or legal action against Prak Vutha.
No legal action was taken against two policemen from Border Protection Unit 701 implicated in a February 2006 beating of a 13‑year‑old boy. Likewise, there was no action against officials and no progress in the police investigation of an April 2006 case involving Police Commissioner Team Sangkriem in Preah Vihear Province and three other police agents who detained Kong Salath without a warrant and beat him. No disciplinary or legal action was taken against abusive officers in the April 2006 beating of a motorist by Battambang military police. Regarding the December 2006 case of Tous Sdoeung, whom two military police officers allegedly tortured to death while in detention, early in the year a provincial court prosecutor completed an investigation and forwarded it to an investigating judge. The court investigation continued. The two alleged perpetrators continued to work in their positions as military police officers.
Prison and Detention Center Conditions
Prison conditions did not meet international standards. Conditions remained harsh and at times were life threatening. Government efforts to improve them continued to be hampered by a lack of funds and weak enforcement. Human rights organizations cited a number of serious problems, including overcrowding, medical and sanitation problems, food and water shortages, malnutrition, and poor security. According to LICADHO, the 18 prisons they monitored had a designed capacity of approximately 6,440 inmates but held a total 9,582 inmates.
There were reports at some prisons that cells of 40 by 20 feet held up to 110 prisoners. At CC1 prison, cells of 26 by 26 feet held an average of 50 prisoners. In some prisons authorities used shackles and held prisoners in small, dark cells as a form of harsher punishment.
LICADHO reported that 56 prisoners in 18 of the country's prisons died during the year.
Government ration allowances for purchasing prisoners' food routinely were misappropriated and remained inadequate, exacerbating malnutrition and disease. One NGO claimed that in some cases prison authorities sold the NGO's donations of supplemental food intended for prisoners. According to rights organizations, families had to bribe prison officials in order to visit prisoners or provide them food and other necessities. NGOs reported that prisoners whose families bribed prison authorities received preferential treatment including access to visitors, transfer to better cells, and the opportunity to leave cells during the day.
There were reports that officials demanded bribes before allowing prisoners to attend trials or appeal hearings and before releasing inmates who had served full jail terms.
In most prisons there was no separation of adult and juvenile prisoners, of male and female prisoners, or of persons convicted of serious crimes and persons detained for minor offenses. Pretrial detainees were routinely held together with convicted prisoners. LICADHO reported that there were 622 incarcerated minors ages 13 to 17, many of whom were held in prisons that did not have facilities to separate minors from adult prisoners.
The government generally continued to allow international and domestic human rights groups, including the International Committee of the Red Cross, to visit prisons and provide human rights training to prison guards. However, NGOs reported that at times cooperation from local authorities was limited. Authorities curtailed access to pretrial detainees, in particular. The MOI continued to require that lawyers, human rights monitors, and other visitors obtain permission prior to visiting prisoners. The MOI withheld such permission in some politically sensitive cases. NGOs were not allowed to interview prisoners in private.
d. Arbitrary Arrest or Detention
The law prohibits arbitrary arrest and detention; however, at times the government did not respect these prohibitions. On June 7, the National Assembly passed a criminal procedures code, and in August the king signed the law into effect. The new code went allows for pretrial detention of up to six months for misdemeanors and 18 months for felonies. Prior to enactment of this code, the maximum length of pretrial detention for an adult person was six months under the UN Transitional Authority in Cambodia (UNTAC) code, although the government sometimes held pretrial detainees for longer periods. ADHOC reported that at least 100 persons were illegally arrested and detained during the year. ADHOC stated that 32 of those illegally detained were subsequently freed following detainee complaints, interventions by human rights NGOs, or payment of bribes. ADHOC believed that the actual number of arbitrary arrests and detentions was somewhat higher, because some victims in rural areas did not file complaints due to difficulty in traveling to the NGO's offices or out of fear for their family's security. According to ADHOC, no legal or disciplinary actions were taken against the persons responsible for the illegal actions.
Role of the Police and Security Apparatus
The General Commissariat of the National Police, which is under the supervision of the MOI, manages all civilian police units. The police forces are divided into those who have the authority to make arrests, those without such authority, and the judicial police. Military police are permitted to arrest civilians on military property or when authorized by local governments.
Police officers acted with impunity, and in most cases the government took little or no action. There were reports that police, prosecutors, investigating judges, and presiding judges received bribes from owners of illegal businesses.
The law requires police, prosecutors, and judges to investigate all complaints, including those of police abuses; however, in practice judges and prosecutors rarely conducted an independent investigation prior to a public trial. Presiding judges passed down verdicts based only on written reports from police and witness testimonies. In general police received little professional training. Police who failed to prevent or respond to societal violence were rarely disciplined.
There were no developments in the April 2006 case of an antidrug department and military police officer who shot and injured a well‑known singer, Sovansocheata. No legal action was taken in the April 2006 case of two Brigade 70 military unit officers who shot and injured a person in Phnom Penh. There were no developments in the June 2006 case in which a military officer shot and injured a garment factory worker. In February an investigating judge in Siem Reap Province issued a warrant for the arrest of three police officers who allegedly raped a 12‑year‑old girl in November 2006; however, the suspects remained at large. There were no developments in the pending appeal of the April 2006 acquittal of three judges, two deputy prosecutors, and two court clerks originally convicted, then retried after appeal on finding of a mistrial, on charges of corruption and corruption‑related conspiracy.
Arrest and Detention
The law requires police to obtain a warrant from an investigating judge prior to making an arrest, but police may arrest without a warrant anyone caught in the act of committing a crime. The law allows police to take a person into custody and conduct an investigation for 48 hours, excluding weekends and government holidays, before charges must be filed. In felony cases of exceptional circumstances prescribed by law, police may detain a suspect for an additional 24 hours with the approval of a prosecutor. However, authorities routinely held persons for extended periods before charging them. Many prisoners, particularly those without legal representation, had no opportunity to seek release on bail. Under the new criminal procedures code, accused persons may be arrested and detained for up to 24 hours before meeting with a lawyer, but prisoners routinely were held for several days before gaining access to a lawyer or family members. According to government officials, such prolonged detention largely was a result of the limited capacity of the court system.
LICADHO reported that as of midyear at least 101 pretrial detainees had been detained longer than the six‑month limit. Under the allowable pretrial detainee periods stipulated by the new code, at year's end there were at least 34 such prisoners.
On May 19, two military police officers in Banteay Meanchey Province detained Kim Heang for three days after Kim Heang had a dispute with his neighbor, a regional military official. The two officers made the arrest without a warrant. After an NGO intervened, the officers' commander ordered Kim Heang released. No administrative or legal action was taken against the officers.
On May 25, the Ratanakiri Provincial Court sentenced a 13‑year‑old Jarai ethnic minority youth to eight months and 10 days in prison for stealing brass gongs. The youth was 12 years old when arrested, under the minimum age for imprisonment, but spent more than eight months in pretrial detention. According to ADHOC, on May 25, a prosecutor filed a suit with the appeals court, but at year's end the youth remained in jail.
On August 9, the Phnom Penh Municipal Court convicted six persons and acquitted two charged with planning bombings at the November 2006 Water Festival. Two of the convicted were sentenced in absentia to 12 years in prison. The remaining four received six‑year sentences. Lawyers and NGOs maintained the police did not serve arrest warrants or tell the suspects the charges against them.
e. Denial of Fair Public Trial
The constitution provides for an independent judiciary, but the government did not respect judicial independence. The courts were subject to influence and interference by the executive branch, and there was widespread corruption among judges, prosecutors, and court officials.
The court system consists of lower courts, an appeals court, and the Supreme Court. The constitution also mandates a constitutional council, which is empowered to review the constitutionality of laws, and a supreme council of the magistracy, which appoints, oversees, and disciplines judges. The composition of both councils heavily favored the CPP.
There is a separate military court system, which suffered from deficiencies similar to those of the civilian court system. While civilians may fall under military court jurisdiction in some cases, the legal distinction between the military and civil courts sometimes was ignored in practice. Civilians have been called for interrogation by military courts with no apparent jurisdiction in their cases.
On June 12, the Extraordinary Chambers in the Courts of Cambodia (ECCC) adopted its internal rules to begin prosecuting egregious crimes of the 1975‑79 Khmer Rouge regime. On July 31, the ECCC coinvestigating judges charged Kaing Guek Eav (alias Duch), former Khmer Rouge director of the S‑21 torture prison, or Tuol Sleng, for crimes against humanity and placed him in an ECCC provisional detention center. The ECCC later arrested and detained four more Khmer Rouge leaders and charged them with crimes against humanity and war crimes: Nuon Chea (also known as "Brother Number 2"), Khieu Samphan, Ieng Sary, and Ieng Thirith, who was charged only with crimes against humanity. In August Duch's lawyers filed an appeal against his provisional detention. On December 3, the ECCC pretrial chamber decided unanimously to affirm the detention order and dismiss the appeal.
Trials are public. Juries are not used; the presiding judge possesses the authority to pass a verdict. Defendants have the right to be present and consult with an attorney, confront and question witnesses against them, and present witnesses and evidence on their own behalf. If a defendant cannot afford an attorney, the court is required to provide the defendant with free legal representation; however, the judiciary lacked the resources to provide legal counsel, and most defendants sought assistance from NGOs or went without legal representation. Trials typically were perfunctory, and extensive cross‑examination usually did not take place. Defendants and their attorneys have the right to examine government‑held evidence relevant to their cases; however, at times it was difficult for them to obtain such access, especially if the case was political or involved a high‑ranking government official or well‑connected member of the elite.
Defendants are entitled by law to the presumption of innocence and the right of appeal, but due to pervasive corruption, defendants often were expected to bribe judges to secure a verdict. A citizen's right to appeal sometimes was limited by difficulty in transferring prisoners from provincial prisons to the appeals court in Phnom Penh. Many appeals thus were heard in the absence of the defendant.
A lack of resources, low salaries, and poor training contributed to a high level of corruption and inefficiency in the judicial branch, and the government did not ensure due process. From January through September, the Center for Social Development monitored 1,420 felony and misdemeanor hearings with 2,437 defendants and found trial procedure abuses in the Supreme Court, appeals court, and four lower courts. In a report of trials observed from January to March, the center stated that courts tried 34 percent of 740 defendants in absentia. At the appeals level, defendants were not present during trial in 69 percent of cases. Of defendants charged with felonies, 37 percent had legal representation, compared with 7 percent of those charged with misdemeanors.
Officials reported many suits were pending due to a shortage of judges and courtrooms. NGOs blamed the slow process on court officials who focused on cases from which they could gain financial benefits.
There remained a critical shortage of trained lawyers, particularly outside Phnom Penh. Persons without means to secure counsel often were effectively denied the right to a fair trial. According to the Bar Association, approximately 30 percent of the country's 573 lawyers provided legal counsel to poor persons, although this was inadequate to cover the basic legal rights of all of the country's poor.
Sworn written statements from witnesses and the accused usually constituted the only evidence presented at trials. The accused person's statements sometimes were coerced through beatings or threats, and illiterate defendants often were not informed of the content of written confessions that they were forced to sign. In cases involving military personnel, military officers often exerted pressure on judges of civilian criminal courts to have the defendants released without trial.
Court delays or corrupt practices often allowed accused persons to escape prosecution. Government officials or members of their families who committed crimes often enjoyed impunity.
Although the courts prosecuted some members of the security forces for human rights abuses, impunity for most of those who committed abuses remained a problem. In many criminal cases, rich or powerful accused individuals usually paid money to victims and authorities to drop the criminal charges against them. Authorities were known to urge victims or their families to accept financial restitution in exchange for dropping criminal charges.
Born Samnang and Sok Sam Oeun remained in prison for the murder of Chea Vichea. On April 12, an appeals court hearing upheld the Phnom Penh Municipal Court decision sentencing them to 20 years each in prison, despite new exculpatory evidence. On June 7, lawyers filed grievances with the Supreme Court and at year's end were awaiting action.
On June 7, the National Assembly passed a criminal procedures code, and in August the king signed the law into effect. As a cornerstone of national law also to be employed at the Khmer Rouge Tribunal, the new code was based on wide international consultation and was viewed as meeting an international standard suitable for the tribunal's trial judges.
Political Prisoners and Detainees
There were no reports of political prisoners or detainees.
Civil Judicial Procedures and Remedies
The country has a judiciary in civil matters, and citizens are entitled to bring lawsuits seeking damages for human rights violations. Generally, there are both administrative and judicial remedies. However, the judiciary was generally viewed as corrupt, politically biased, and weak, and persons seldom filed complaints because they did not trust the judicial system. The public appeared especially distrusting of the judiciary to act in a transparent manner when a case was in conflict with the government. Enforcing a court order for a civil or criminal case was often problematic. Persons occasionally turned to vigilante justice.
f. Arbitrary Interference with Privacy, Family, Home, or Correspondence
The law provides for the privacy of residence and correspondence and prohibits illegal searches; however, police routinely conducted searches and seizures without warrants.
There continued to be reports of authorities entering private properties without proper judicial authorization. Due to the forced collectivization during Khmer Rouge rule and the return of thousands of refugees, land ownership often was unclear, and most landowners lacked adequate formal documentation of ownership. The 2001 land law states that any person who peacefully possessed private property without contention for five years prior to the 2001 promulgation of the law has the right to a definitive title to that property. Widespread land speculation fueled disputes and increased tensions between poor rural communities and speculators. The Cadastral Commission, which settles disputes over land that was not registered or where an owner was not given a land certificate, continued to perform its functions slowly. The courts remained responsible for resolving disputes in cases where land was registered or disputants were given land titles. The National Authority for Resolving Land Disputes, established in 2006 to adjudicate land cases, was ineffective.
Problems of inhabitants being forced to relocate continued to occur when officials or businesspersons colluded with local authorities. Some persons also used the court system to intimidate the poor and vulnerable into exchanging their land for compensation below market value. ADHOC reported receiving 382 land‑related cases affecting 19,329 families during the year. During the same period, LICADHO received 98 land‑related cases in Phnom Penh and 13 other provinces affecting a total of 6,048 persons. The poor often had no legal documents to support their land claims and lacked faith in the judicial system. Some of those expelled successfully contested these actions in court, but the majority lost their cases.
On January 23, 200 ethnic Jarai villagers in Ratanakiri Province filed a complaint with the provincial court and a criminal complaint with the provincial prosecutor against Keat Kolney, a well‑connected individual, for confiscation of 1,112 acres of their land in 2004. Many villagers rejected Keat Kolney's settlement offers. On June 19, Keat Kolney sent a letter to the Cambodian Bar Association alleging the legal aid NGO lawyers who represented the villagers trained the villagers to say false things to the media and asked the association to investigate the lawyers. On June 21, Keat Kolney filed criminal complaints accusing the villagers of fraud and the lawyers of inciting villagers to commit fraud. In late July 42 of the 200 villagers retracted earlier statements and said they willingly sold the land to Keat Kolney. At year's end a provincial court investigation continued.
On April 20, approximately 150 military police and police officers armed with guns, electric batons, and tear gas forcibly evicted 117 families from the Mittapheap District in Sihanoukville. Several villagers were injured, and their houses were demolished.
On May 4, 30 L'vea Em District villagers of Kandal Province approached the National Assembly to intervene in an economic land concession. Community families had been farming and inhabiting the disputed 1,730 acres when, on April 25, a Chinese company began digging up the land, acting on a 2006 government concession for development. The company reportedly suspended operations due to villagers' protests.
On May 29, 40 soldiers from ACO Tank Command Headquarters in Kampong Speu Province used an armored vehicle to destroy crops and fences on 60 acres of land occupied by 25 Phnom Srouch District families. Military officials stated the land was part of a shooting range and the villagers had illegally occupied the land. In 2002 the villagers had sought title to their land with the Cadastral Commission, and in 2006 they complained to the National Authority to Resolve Land Disputes, stating they had lived on the land since 1979. At year's end these requests had not resulted in any action.
Eviction notices were served without proper judicial authorization. On July 31, Sihanoukville City Hall issued an eviction notice ordering more than 100 families out of the city's Mittapheap District. Responding to villagers' plea for intervention, the prime minister ordered the Sihanoukville governor to reexamine the case. On May 8, representatives of 146 families of the Phnom Penh Tonle Bassac Group 78 (G78) area made public their own neighborhood development plan. The plan was in response to a June 2006 municipality eviction notice stating the land would be developed for beautification purposes. Many of the families had lived on the land since the 1980s and claimed ownership under the 2001 land law. G78 community members stated that the municipality offered compensation that was approximately equal to 2 percent of the independently assessed market value, plus one plot in a Phnom Penh eviction resettlement site per family. At year's end there were no decisions on these evictions.
On March 3, the CPP Central Committee granted the prime minister sole power to resolve land disputes involving CPP officials. The prime minister then announced a war against illegal land grabbers, warning CPP officials to surrender illegally obtained land or face removal from their positions. As a result, the government claimed that the director general of the military's technical and materials department, Chao Phirun, handed over 495 acres of land and an anonymous CPP official withdrew from a Kandal Province land dispute. Neither official faced reprimand. On March 10, authorities arrested CPP member and military Colonel Te Haing over a 2,500‑acre land dispute in Banteay Meanchey Province. At year's end Te Haing was awaiting trial. In November Tan Seng Hak, a former advisor to CPP Chairperson Chea Sim, was convicted for falsifying documents and giving false testimony in connection with his alleged efforts to take over 740 acres in Phnom Penh. He was sentenced to a total of five years and eight months in prison. There were no developments on the prime minister's May 2006 plan to redistribute 494,210 acres of land to 50,000 farmers in Sihanoukville.
Living conditions worsened at two of the resettlement sites for former residents of the two Phnom Penh communities of Tonle Bassac Sambok Chab and Preah Monivong Hospital areas, whose 1,200 and 168 families, respectively, were evicted in 2006, reportedly through a nontransparent process that may not have included proper judicial authorization. Authorities offered evicted residents relocation site plots, but plots at two of the sites were widely considered to be inadequate compensation. The sites lacked sufficient sanitation facilities, electricity, clean water, health facilities, schools, and central markets and were far from Phnom Penh's commercial center, where residents could earn an adequate income.
The appeals court took no action on a February 2006 complaint by SRP parliamentarian Son Chhay, who was directed by the Siem Reap Provincial Court to sell 7.8 acres of his land to a government agency for an amount below the market price. The appeals court took no action in the 2006 case of 12 persons convicted in connection with a Kampot Province confrontation between 2,000 squatters and local police over the squatters' rights to live on government land. In a 2006 eviction case in Peam Chor District, Prey Veng Province, that left one person dead and four others injured, police who were implicated in the killing accused some of the villagers of robbery, in what NGOs said was an attempt to intimidate the villagers. On November 1, the provincial court questioned seven of the villagers on the robbery charges, and the court investigation continued at year's end. However, there were no new developments in the investigation of the eviction killings and injuries. There were no new developments in the August 2006 Koh Kong Province land dispute in which the Ministry of Agriculture provided two adjacent land concessions to businessperson and CPP senator Ly Yong Phat in contravention of known legal standards.
There were no developments in a 2005 land dispute involving indigenous Phnong hill tribe members and a Chinese company in Mondulkiri Province.
Section 2 Respect for Civil Liberties, Including:
a. Freedom of Speech and Press
The constitution provides for freedom of speech and of the press; however, these rights were not always respected in practice.
The constitution implicitly limits free speech by requiring that it not adversely affect public security. The constitution also declares that the king is "inviolable." In December the Ministry of Information issued a directive that reiterates these limits and prohibits publishers and editors from running stories that insult or defame government leaders and institutions.
The 1995 press law prohibits prepublication censorship or imprisonment for expressing opinions. However, the government continued to use the older UNTAC law to prosecute journalists and others on defamation and disinformation charges. In 2006 the National Assembly amended the UNTAC law, eliminating imprisonment for defamation but not for spreading disinformation, which carries prison sentences of up to three years. In both types of cases, judges can order fines, which may lead to jail time if not paid.
The government and influential individuals used the weak and often politically biased judiciary to file defamation and disinformation suits, both civil and criminal, in an effort to silence critics. In February the Phnom Penh Municipal Court charged Sralanh Khmer of disinformation and insulting the court's director, Chiv Keng. Also in February Sihanoukville Municipal Governor Say Hak filed a defamation suit against pro‑Norodom Ranariddh Party (NRP) newspaper Samleng Yuveachun Khmer (Voice of Khmer Youth) over an article linking him to land grabbing. In July NRP Vice Secretary General Sao Rany filed a defamation suit against Sralahn Khmer for printing a report claiming his daughter had an affair with Prince Ranariddh. In November General Un Den lodged a defamation and disinformation lawsuit against Thach Keth, the publisher of Sralanh Khmer, for printing an article that accused the general of smuggling vehicles across the border from Thailand. At year's end there were no formal decisions in these cases.
In July Phnom Penh Municipality Governor Kep Chuktema filed a disinformation suit against the editor of Samleng Yuveachun Khmer for an article alleging the governor sold City Hall to a private developer. In November the editor paid bail of $500 (two million riel), and at year's end the case was pending with the court.
The constitution states that the country shall not invade any country nor interfere in any other country's internal affairs, directly or indirectly. Making a statement in contravention of this constitutional provision is considered a crime. In the case of at least one Khmer Kampuchea Krom activist, an arrest warrant was outstanding due to his statements about what the government considers to be sovereign Vietnamese territory.
Many interpreted a law passed in 2006 as limiting the right of members of Parliament (MPs) to speak freely. The law declares that MPs may not use their parliamentary immunity to abuse national security, social customs, or an individual's honor. In addition, the law allows an MP to be arrested, charged, and detained prior to the lifting of parliamentary immunity. At year's end no MP had been charged under this law.
All major political parties had reasonable and regular access to the print media. Although the press law does not specifically permit newspapers to receive financial support from political parties, all major Khmer‑language newspapers received such support and were politically aligned. There were an estimated 20 Khmer‑language newspapers published regularly; more than half were considered pro‑CPP, and at least two newspapers were considered to support each of the other main political parties--FUNCINPEC, the SRP, and the NRP. Although the three largest circulation newspapers were considered pro‑CPP, most newspapers criticized the government, particularly on corruption and land grabbing. The prime minister, NRP president Prince Norodom Ranariddh, FUNCINPEC party leaders, and opposition party leader Sam Rainsy frequently came under attack.
The government, military forces, and ruling political party continued to dominate the broadcast media and influence the content of broadcasts. There were approximately 50 radio stations and seven television stations. Most were controlled or strongly influenced by the CPP, although a few were independent or aligned with other parties. In August the Ministry of Information issued a broadcast license to a CPP government official to open a Phnom Penh radio station after denying similar requests since 2003 from the SRP and independent human rights advocacy organizations. In September the Cambodian Center for Human Rights (CCHR) Voice of Democracy (VOD) again requested a license, but the ministry denied the application, restating previous claims that the media market was saturated. In February the Information Ministry ordered all television and radio stations not to transfer or sell licenses if unable to continue operating and to return the license to the ministry.
Despite being unable to obtain a broadcast license, the VOD radio program, which included independent and often antigovernment views, remained popular and continued broadcasting its program on several radio stations, including the SRP‑aligned radio station FM 93.5. In July CCHR announced that it had transferred management of VOD to the newly established Cambodian Center of Independent Media. Taped programming from Voice of America (VOA) and Radio Free Asia (RFA) Khmer‑language service was also regularly broadcast on Beehive/FM 105, the Women's Media Center FM 102, and Rota Angkor FM 95.5 (Siem Reap) radio stations. Four political parties were each broadcasting daily one‑hour shows on FM 105.
Journalists, publishers, and distributors were also subject to other forms of harassment and intimidation, including death threats. In June the Ministry of Information ordered the confiscation of printed copies of a report by the international NGO Global Witness accusing the prime minister and close relatives and associates of involvement in illegal logging. The report was freely available via the Internet, and local media made references to the report. A June 8 letter signed by Information Minister Khieu Kanharith to Sralanh Khmer demanded the newspaper immediately stop publishing a serialized version of the report or face legal action; the newspaper complied.
In the same period, French‑language newspaper Cambodge Soir closed down, reportedly due to bankruptcy. A few days prior, according to media reports, employees declared a strike because a reporter was dismissed for publishing a story about the Global Witness report. After a closure of several months, the newspaper resumed as a weekly publication.
In October Radio Beehive Director Mam Sonando suspended the NRP's radio show "Royalist Voices" for several days after the program criticized the prime minister. An RFA journalist fled the country reportedly after receiving death threats for his coverage of illegal logging in Kampong Thom Province. The reporter returned several weeks later and resumed work unharmed.
In February the Ministry of Information threatened to close pro‑NRP newspaper Khmer Amatak (Permanent Khmer) for printing an article alleged to have incited tension between the country's two main Buddhist sects, and in October the ministry suspended the newspaper for one month for failing to publish a "correction" the ministry requested regarding a September article involving two senior FUNCINPEC officials.
In early November authorities seized copies of the premier issue of foreign‑funded Free Press Magazine for criticizing retired king Norodom Sihanouk, the prime minister, and other government officials.
In May the prime minister publicly criticized an RFA reporter as "insolent" and "rude" for asking questions about the coalition between the CPP and FUNCINPEC. Purportedly fearing for his personal safety, the reporter went into hiding outside the country. He returned to work a few weeks later without incident.
In June three reporters from the newspapers Sarpotamean Ekkereach Kampuchea (Cambodian Independence News) and Sarpotamean Tasanak Khmer (Khmer Vision News) in Pursat Province alleged that provincial court official bodyguards beat them at gunpoint for trying to take photographs of a truck carrying illegal timber.
On August 4, Oeun Vannak, deputy commander of the Pursat Province military police, allegedly physically attacked journalist Heng Veasna over the journalist's investigation into claims of illegal use of firearms by two military police officers. Heng Veasna filed a complaint with the provincial prosecutor, but at year's end the court had not taken action in the case, and Oeun Vannak continued his military police duty.
In August the home of a Chhbas Ka (Accurate News) newspaper reporter was set on fire twice. The journalist claimed the first fire was set the day after he received a threatening telephone call over his report on illegal logging in Pursat Province. A few days later, provincial authorities charged two suspects with arson. The cases were pending before the court at year's end.
In November a man accused three journalists from the newspaper Samleng Santepheap (Voice of Peace) in Kampong Thom of stealing $1,050 (4.2 million riel) when the reporters visited his home to investigate allegations that he was illegally raising snakehead fish. Police arrested and questioned the reporters but eventually released them after the intervention of a senior official from the Ministry of Information.
In December a VOD reporter investigating the removal of a statue from a pagoda claimed to have been detained by military and police officers, who deleted photographs from his camera before releasing him.
Most reporters and editors privately admitted to some self‑censorship due to fear of government reprisals. In February two major daily Khmer newspapers refused to print advertisements demanding justice for the two men imprisoned for the killing of union leader Chea Vichea. Reporters for VOA, RFA, and some opposition newspapers worked from unmarked offices and reported stories using pseudonyms.
The government‑controlled national television and radio stations broadcast live or taped sessions of National Assembly debates; however, in several instances these broadcasts were censored. National radio and television stations broadcast some human rights, social action, public health, education, and civil society programming produced by domestic NGOs.
The government occasionally restricted media access to some government facilities. The constitution mandates media access to National Assembly sessions, and the National Assembly allowed reporters to enter its grounds upon clearance by its security office. In 2005 the Phnom Penh Municipal Court chief ordered that reporters must have written permission to bring recording devices into the courtroom and to interview court officials. Such permission rarely was sought, and there were no reports of the court denying permission. A July 2006 Council of Ministers directive prohibiting government officials and employees from speaking to the media or the public about government corruption remained in effect.
In August government authorities confiscated digital video discs with images of bodies in an airplane crash in Kampot Province, stating that the images would create fear among tourists. The video discs remained readily available in Phnom Penh and other areas of the country.
There were no government restrictions on access to the Internet or reports that the government monitored e‑mail or Internet chat rooms. Individuals and groups could engage in the peaceful expression of views via the Internet, including by e‑mail. Although the International Telecommunication Union estimated the country's Internet penetration was 0.3 percent in 2006, in urban areas Internet access was widely available through Internet cafes and home subscriptions.
Academic Freedom and Cultural Events
In general there were no legal impediments to academic freedom. However, scholars tended to be careful when teaching politically related subjects for fear of offending politicians. In February the Phnom Penh Municipal Court sentenced Tieng Narith, a former professor at the Buddhist University of the Royal Academy of Preah Sihanouk Reach, to two and a half years in prison for teaching from a self‑published text containing antigovernment material. The verdict also ordered a fine of $1,250 (5.25 million riels) or two additional years in prison. Tieng Narith's family claimed that he was mentally ill, and during the trial the court ordered a psychiatric examination, the results of which were kept confidential. It was unclear how the medical exam results affected the case, if at all. At year's end the case was under appeal.
There were no government restrictions on cultural events.
b. Freedom of Peaceful Assembly and Association
Freedom of Assembly
The constitution provides for freedom of peaceful assembly, but at times the government did not respect this right in practice. The government required that a permit be obtained in advance of a march or demonstration. The government routinely did not issue permits to groups critical of the ruling party or of nations with which the government had friendly relations. Authorities cited the need for stability and public security as reasons for denying permits. Police forcibly dispersed groups that assembled without a permit, often resulting in minor injuries to some demonstrators.
ADHOC reported that out of 98 protests--55 of which were related to land and 26 to labor disputes--police and military police dispersed 17 protests, three of which were by labor protesters, 10 by land rights protesters, and four by Khmer Kampuchea Krom monks. However, the government permitted some human rights‑related marches and demonstrations. On December 10, the government permitted a Human Rights Day march of 500 human rights activists, monks, and other persons and a rally of an estimated 2,500 persons in Phnom Penh. In the previous two years, such rallies without marches occurred in an enclosed space.
On February 27, police and military police dispersed 60 Khmer Kampuchea Krom Buddhist monks demonstrating at the Vietnamese embassy in Phnom Penh during a state visit by the Vietnamese president. Demonstrators assembled to support Khmer Kampuchea Krom monks in Vietnam who had been defrocked and arrested, urging their release and reinstatement as monks. The next morning one monk protester was found dead with his throat cut. On March 16, police and local authorities in Kandal Province prevented the deceased monk's Khmer Kampuchea Krom community members and monks from holding his funeral.
On April 20, police and municipal authorities dispersed 80 Khmer Kampuchea Krom monks assembled at the Vietnamese embassy trying to deliver a petition in protest of alleged Vietnamese government rights abuses of Khmer Kampuchea Krom living in Vietnam. The protesters decided to go to another embassy to present the petition. On the way a group of unidentified, non‑Khmer Kampuchea Krom monks and laypersons aggressively intercepted the demonstrators and attempted to disperse them. In the ensuing scuffle, one of the Khmer Kampuchea Krom monks was injured. Authorities did not intervene in the confrontation and did not conduct an investigation. On December 17, 40 monks sought again to deliver a petition to Vietnamese embassy officials for the release of Tim Sakhorn and other Khmer Kampuchea Krom monks imprisoned in Vietnam, and also for the return of land that they claimed the Vietnamese government seized from Khmer Kampuchea Krom persons in southern Vietnam. Police attempted to disperse the crowd, but the monks refused to disband, and violence broke out on part of both the police and the monks. A local NGO reported that six monks were injured; police stated that some of the police sustained minor injuries.
On November 26, Ratanakiri provincial police blocked the CCHR from holding a public forum in Kong Yu Village, where community members were embroiled in a land dispute with Keat Kolney. Police gave conflicting reasons for preventing the forum from taking place (see section 1.f.).
On June 8, Supreme Patriarch Non Ngeth and Minister of Cults and Religious Affairs Khun Haing signed a directive prohibiting monks from participating in protests, strikes, riots, or marches. According to media reports, a constitutional council member stated the ban violated the constitution.
Freedom of Association
The constitution provides for freedom of association, and the government generally respected this right in practice; however, the government did not effectively enforce the freedom of association provisions of the labor law.
Membership in the Khmer Rouge, which ruled the country from 1975 to 1979 and after its overthrow conducted an armed insurgency against the government, is illegal, as is membership in an armed group.
c. Freedom of Religion
The constitution provides for freedom of religion, and the government generally respected this right in practice. The constitution also prohibits discrimination based on religion, and minority religions experienced little or no official discrimination. Buddhism is the state religion, and more than 93 percent of the population is Buddhist. Ethnic Cham Muslims constitute most of the remaining population.
The law requires all religious groups, including Buddhists, to submit applications to the Ministry of Cults and Religious Affairs to construct places of worship and conduct religious activities. However, there is no penalty for failing to register. In July the Ministry of Cults and Religious Affairs issued a directive restating a 2003 order prohibiting public proselytizing, which continued to be loosely enforced. On August 10, authorities in Phnom Penh dispersed a gathering of approximately 3,000 Christians, stating that organizers did not have proper permits. Prior to the gathering, organizers obtained a permit from the MOI but had not received a response on a request pending with the Ministry of Cults and Religious Affairs.
Societal Abuses and Discrimination
Minority religions experienced little or no societal discrimination. There was no known Jewish community in the country, and there were no reports of anti‑Semitic acts.
For a more detailed discussion, see the 2007 International Religious Freedom Report.
d. Freedom of Movement, Internally Displaced Persons, Protection of Refugees, and Stateless Persons
The law provides for freedom of movement within the country, foreign travel, emigration, and repatriation, and the government generally respected these rights in practice.
The constitution prohibits forced exile, and the government did not employ it.
Protection of Refugees
The laws provide for the granting of asylum or refugee status in accordance with the 1951 UN Convention relating to the Status of Refugees and its 1967 protocol, and the government has established a system for providing protection to refugees. The convention and its protocol have had the full force of the law in the country since accession in 1992. According to the Office of the UN High Commissioner for Refugees (UNHCR), the government abides by the convention and international customary law on refugees. The government allows the UNHCR to process asylum seekers and assist refugees while they are in the country.
A memorandum of understanding that the government signed in 2005 with the UNHCR and the government of Vietnam to resolve the situation of Montagnards under UNHCR protection remained in effect. Asylum seekers who reached the UNHCR Phnom Penh office were processed with government cooperation. During the year there were 449 new arrivals seeking asylum with the UNHCR. According to the UNHCR, 97 Montagnard and 20 non‑Montagnard refugees departed for a third country, while authorities deported 30 rejected Montagnard asylum seekers to Vietnam, and 33 Montagnards voluntarily returned to their country of origin. At year's end there were 467 Montagnards in UNHCR protection sites in Phnom Penh, which included 101 Montagnards who arrived in previous years. According to the UNHCR, during the year no refugees requested local integration.
In practice the government provided some protection against refoulement, the return of persons to a country where there is reason to believe they feared persecution. Through the assistance of the UNHCR, during the year the government provided temporary protection to individuals who may not qualify as refugees under the 1951 convention and the 1967 protocol, affording such protection to approximately 150 persons. However, an NGO based in Ratanakiri Province reported that local police unofficially returned 59 asylum-seeking Montagnards to Vietnam without UNHCR review.
On April 20, Ratanakiri provincial police arrested two persons on charges of human trafficking for their roles assisting Montagnards. NGOs claimed the suspects provided asylum seekers food, shelter, and transportation to the UNHCR office in Phnom Penh once they had crossed the border from Vietnam. At year's end the suspects had been released with charges against them dropped.
An NGO claimed that local authorities at the border with Vietnam continued searches for Montagnards when they received information about new arrivals of Montagnards. There were unconfirmed reports that Vietnamese authorities offered incentive awards to Cambodian border police who returned Vietnamese refugees to Vietnam and that Vietnamese secret police covertly conducted searches for Vietnamese refugees on the Cambodian side of the border.
The country had habitual residents who were de facto stateless, and the government had not effectively implemented laws or policies to provide such persons the opportunity to gain nationality. Under the nationality law, citizenship is derived by birth from a foreign mother and father who were born and living legally in the country, or from a mother or father who has Cambodian citizenship. A study commissioned by the UNHCR estimated that several thousand potentially stateless persons lived in the country. However, the study's estimated number of such persons came from anecdotal evidence of NGOs that provided services to disenfranchised communities, including persons with no proof of nationality, and not from a survey of stateless persons; therefore, local UNHCR representatives did not consider the figure conclusive.
The UNHCR stated that the country's potentially stateless population included mostly ethnic Vietnamese. According to an NGO that worked with ethnic Vietnamese, individuals without proof of nationality often did not have access to formal employment, education, marriage registration, the courts, and land ownership. The most common reason for statelessness was lack of proper documents from the country of origin.
Section 3 Respect for Political Rights: The Right of Citizens to Change Their Government
The constitution provides citizens the right to change their government peacefully, and citizens generally exercised this right in practice through periodic elections on the basis of universal suffrage. Suffrage is voluntary for all citizens age 18 years and older.
Elections and Political Participation
On April 1, the country held elections for 11,353 chiefs, first deputies, second deputies, and councilors for 1,621 commune councils. The CPP won 70.4 percent of the positions, the SRP 23.4 percent, NRP 3.7 percent, and FUNCINPEC 2.4 percent.
Most observers agreed the commune council elections were the least violent and best organized elections ever held in the country. While there were problems at some polling stations, NGOs, opposition parties, monitors, and others disagreed as to how and whether the problems affected the overall outcome of the elections. Three NGOs reported that election officials did not allow some registered voters to vote because of voter registration list discrepancies such as mistyped and misspelled names or absence of names from the voter list, often due to names transferred to different polling stations without informing the voter. Additionally, NGOs and opposition parties complained that the CPP started advertising weeks or months in advance through the mostly CPP‑dominated media without reprisal. NGOs reported that on election day, some ruling party incumbents illegally issued voter registration documents, stood watch in some polling station areas where local authorities were prohibited, and provided assistance to voters in these prohibited areas.
Parties and individuals were free to be candidates without restrictions. On March 13, the Phnom Penh Municipal Court sentenced NRP president Prince Norodom Ranariddh in absentia to 18 months in prison and a $150,000 (600 million riel) fine on charges of breach of trust. The prince chose self‑exile during the election campaign and on election day rather than face the charges. On October 3, the appeals court rejected his appeal of the Phnom Penh court's decision. Plans to appeal this decision had not been realized by year's end.
Some NGOs and political parties alleged that membership in the dominant CPP party provided advantages, such as gifts or access to government emergency aid.
Traditional culture limited the role of women in government; however, women took an active part in the 2003 national elections. The number of women increased in the National Assembly, Senate, and senior government positions. There were 22 women in the 123‑seat National Assembly, nine women in the 61‑seat Senate, and 24 women working as ministers, secretaries of state, undersecretaries of state, and National Election Commission officials. Women also served as advisors, and there were 13 female judges in the municipal and provincial courts, appeals court, and Supreme Court. In the April commune council elections, 14.6 percent of the elected councillors were women, of whom 67 were elected as chiefs. This was an increase from the 2002 commune council elections, when women won 9 percent of the total positions.
There were four members of minorities--two Cham and two other ethnic minorities--in the National Assembly. There also were six members of minorities in the Senate. At least eight officials in senior positions in the government were from minority groups.
Government Corruption and Transparency
There is no anticorruption law, and only a few provisions of other laws provide criminal penalties for official corruption. Officials frequently engaged in corrupt practices with impunity. The World Bank's Worldwide Governance Indicators reflected that corruption was a severe problem.
In 2005 the prime minister instructed the Ministry of National Assembly‑Senate Relations and Inspection to prepare a draft anticorruption law. As of year's end, observers had not seen a revised draft since September 2006, and the issue was pending with the Council of Ministers.
Corruption was considered endemic and extended throughout all segments of society, including the executive, legislative, and judicial branches of government. Public perception of corruption was widespread. A 2006 Economic Institute of Cambodia assessment found that the private sector perceived the judiciary to be the most corrupt institution in the country, followed by the tax and customs services, public health care, and police. Meager salaries contributed to "survival corruption" among low‑level public servants, while a culture of impunity enabled corruption to flourish among senior officials.
The Economic Institute's 2006 assessment of corruption in the private sector estimated that in 2005 private sector unofficial payments to public officials totaled $330 million riels. The assessment also found that the larger the private firm, the larger the payments required by government officials. The same study found that approximately 25 percent of potential taxes were collected from the private sector in 2005, representing a loss to the government of approximately $400 million riels. In June Global Witness published a report charging high‑level government officials with corruption related to illegal logging. Some observers and many government officials criticized the report as noncredible based on its heavy reliance on anonymous sources.
The National Archives Law allows unlimited access to informational documents in the public archive. However, the law grants access to other unspecified government documents only after 20 years, and documents affecting national security and preservation of personal lives would be released after 40 and 120 years, respectively. In practice the government occasionally denied access to information, citing reasons of confidentiality or national security.
Section 4 Governmental Attitude Regarding International and Nongovernmental Investigation of Alleged Violations of Human Rights
A variety of domestic and international human rights groups generally operated without government restriction, investigating and publishing their findings on human rights cases. Government officials often cooperated with human rights workers in performing their investigations; however, there were numerous reports of lack of cooperation or even intimidation by local authorities throughout the country.
There were approximately 40 human rights NGOs in the country, but only a small portion of them were actively involved in organizing training programs or investigating abuses.
Domestic and international human rights organizations faced threats and harassment from local officials. These took the form of restrictions on and disruptions during gatherings sponsored by NGOs, verbal intimidation, threats of legal action, and bureaucratic obstruction.
On May 15, a CCHR coordinator went into hiding claiming that Sihanoukville authorities threatened to arrest him on charges of forming an illegal armed force for his role organizing a resistance effort to a forced eviction in Sihanoukville's Mittapheap District. He returned to work a few weeks later without incident.
There were no developments in the May 2006 case of an ADHOC activist temporarily detained in Koh Kong Province for photographing a confrontation between villagers and officials.
In January the UN Special Representative of the Secretary‑General for human rights in Cambodia, Yash Ghai, submitted to the UN Human Rights Council a report on his March 2006 visit expressing concerns about land grabbing and government land concessions. Afterward the prime minister publicly called the special representative derogatory names, refusing to meet with him ever again. In December Yash Ghai made a 10‑day assessment visit to the country during which the prime minister reiterated his opinions and no government official granted him a meeting.
The Cambodian National Human Rights Commission remained largely inactive. The committee did not have regular meetings or a transparent operating process.
Section 5 Discrimination, Societal Abuses, and Trafficking in Persons
The constitution prohibits discrimination based on race, sex, color, or language; however, the government did not generally protect these rights.
The law prohibits rape and assault; nevertheless, local and international NGOs reported that violence against women, including domestic violence and rape, was common. Rape is a criminal offense and punishable by a prison sentence of between five and 10 years, according to the UNTAC law. Spousal rape and domestic abuse are not recognized as separate crimes. A case of spousal rape could be prosecuted as "rape," "causing injury," or "indecent assault," but such charges were rare. The domestic violence law criminalizes domestic violence but does not specifically set out penalties. However, the UNTAC law on battery and injury can be used to penalize domestic violence offenses, with penalties ranging from two months to five years' imprisonment.
According to one NGO, there were 1,025 cases of domestic violence and 221 cases of rape reported in three provinces. Of these cases, courts tried 104 and 28, respectively, resulting in successful conviction in five cases of domestic violence and 20 cases of rape. LICADHO documented 209 cases of domestic violence affecting 213 victims in 12 provinces during the same time period. The MOI's antitrafficking department investigated 529 cases of violence against women and children, resulting in the arrest of 582 perpetrators and rescue of 771 victims. Of the 582 arrests, 458 were for rape and attempted rape. Twelve cases of rape resulted in the death of the victims. A legal advocacy NGO reported receiving 84 cases of domestic violence, 34 of which went to trial during the year. The number of cases likely underreported the scope of the problem, due to ineffective enforcement and the fact that women were afraid to make complaints against perpetrators. NGOs reported that enforcement of the domestic violence law was weak, authorities continued to avoid involvement in domestic disputes, and victims frequently were reluctant to pursue formal complaints.
The government supported NGOs that provided training for poor women vulnerable to spousal abuse, prostitution, and trafficking. A local media center, an NGO, and the Ministry of Women's Affairs produced programming on women's issues. NGOs provided shelters for women in crisis.
The constitution prohibits prostitution; however, there is no specific legislation against working as a prostitute. Trafficking in women for the purpose of prostitution was a serious problem, despite laws against procuring and kidnapping for purposes of sexual exploitation. There were reports that police abused prostitutes. Despite sporadic crackdowns on brothel operators in Phnom Penh, prostitution and related trafficking persisted. Estimates of the number of working prostitutes ranged from 14,725 to 18,250. Sex tourism was a problem, fueled by pervasive poverty and the perception of impunity.
The labor law has provisions against sexual harassment in the workplace but does not specify penalties. The International Labor Organization (ILO) reported that sexual harassment in the industrial sector was rare.
The constitution contains explicit language providing for equal rights for women, equal pay for equal work, and equal status in marriage. In practice women had equal property rights, the same legal status to bring divorce proceedings, and equal access to education and some jobs; however, cultural traditions continued to limit the ability of women to reach senior positions in business and other areas. Women often were concentrated in low‑paying jobs and largely were excluded from management positions. Men made up the vast majority of the military, police, and civil service.
The Ministry of Women's Affairs, mandated to protect the rights of women and promote gender equality in society, continued its Neary Ratanak (Women as Precious Gems) program. The program aimed to improve the image of women through gender mainstreaming, enhanced participation of women in economic and political life, and protection of women's rights.
The constitution provides for children's rights, and the government made the welfare of children a specific goal. The government relied on international aid to fund most child social welfare programs, resulting in only modest funds for problems that affect children.
In 2002 the government instituted a modernized birth registration system administered by the MOI, which reported the program successfully registered 91 percent of births in 2006. The system did not include special outreach to minority communities. The government failed to register all births, resulting in discrimination, including the denial of public services. A study commissioned by the UNHCR on statelessness in the country stated that the birth registration process often excludes children of ethnic minorities and stateless persons. NGOs that provided services to disenfranchised communities reported that children without birth registration were often denied access to education and healthcare. They stated that later in life the same individuals may be unable to access employment, own property, vote, and use the legal system.
Children were affected adversely by an inadequate educational system. Education was free, but not compulsory, through grade nine. Many children left school to help their families in subsistence agriculture, began school at a late age, or did not attend school at all. In 2005 the Ministry of Education reported that 91 percent of eligible children were enrolled in primary school, but this number did not reflect attendance. After primary school, 26 percent of eligible students attended junior high and 9 percent attended high school. Despite a school construction program, schools were overcrowded and lacked sufficient equipment. In rural areas schools often provided only a few years of education. According to ministry data, 46 percent of schools lacked drinking water and 37 percent had no toilets. Teacher salaries were irregularly paid and inadequate to support a decent standard of living, leading to demands for unofficial payments from parents, which poor families could not afford. The government did not deny girls equal access to education; however, families with limited resources often gave priority to boys. In many areas schools were remote and transportation was a problem. This especially affected girls due to safety concerns in traveling between their homes and schools.
Boys and girls had equal access to state‑provided medical care.
Child abuse was believed to be common, although statistics were not available. Child rape remained a serious problem; a local NGO reported 199 cases of rape and attempted rape committed on persons under age 18, two of which resulted in death. Twenty‑nine of the cases involved children below age five. Sexual intercourse with a person under age 15 is illegal; however, child prostitution and trafficking in children occurred. During the year raids on brothels rescued underage girls trafficked for prostitution. The MOI reported arrests of seven foreign pedophiles. Some children were engaged in prostitution for survival, without third‑party involvement.
A domestic NGO estimated that more than 1,200 street children in Phnom Penh had no relationship with their families and more than 10,000 children worked on the streets but returned to families in the evenings. An estimated 500 to 1,500 children lived with their families on the streets in provincial towns. The Ministry of Social Affairs and Youth Rehabilitations provided lower statistics, reporting 3,084 street children nationwide in 2005.
A study conducted by a local NGO stated that in September 2006 there were 37 children under the age of six living with their mothers in prison, and those children were subjected to mistreatment by prison guards and faced physical dangers from adult criminal cellmates. The children generally lacked proper nutrition and education.
Child labor was a problem in the informal sector of the economy.
Trafficking in Persons
The law prohibits trafficking in persons; however, the country was a source, destination, and transit country for men, women, and children trafficked for sexual exploitation and labor. A 2003 study estimated the number of trafficking victims in the sex industry to be 2,000, approximately 80 percent of whom were ethnic Vietnamese women and girls. Children were trafficked domestically for sexual exploitation and labor. Some Vietnamese women and girls were trafficked through the country for exploitation in the commercial sex trade in other Asian countries.
Children were trafficked to Thailand and Vietnam for begging, soliciting, street vending, and flower selling. The children frequently were placed into debt bondage to beg or sell, or they formed part of organized begging rings even when there was no debt or economic hardship involved. Women as well as children were trafficked to Malaysia and Thailand for sexual exploitation and forced labor in factories or as domestic servants, while men were trafficked for forced labor in the agriculture, fishing, and construction sectors.
Trafficking victims, especially those trafficked for sexual exploitation, faced the risk of contracting sexually transmitted diseases, including HIV/AIDS. In some cases victims were detained and physically and mentally abused by traffickers, brothel owners, and clients.
Local traffickers covered specific small geographic areas and acted as middlemen for larger trafficking networks. Organized crime groups, employment agencies, and marriage brokers were believed to have some degree of involvement. Traffickers used a variety of methods to acquire victims. In many cases victims were lured by promises of legitimate employment or travel documents. In other cases acquaintances, friends, and family members sold the victims or received payment for helping deceive them. Young children, the majority of them girls, were often "pledged" as collateral for loans by desperately poor parents; the children were responsible for repaying the loan and the accumulating interest. A September report by the International Organization for Migration (IOM) stated that child domestic workers, particularly those used as collateral or placed into debt bondage, were more likely to be trafficked and to enter commercial sexually exploitive activities.
The law establishes a prison sentence of 15 to 20 years for a person convicted of trafficking in persons under 15 years of age; the penalty is 10 to 15 years for trafficking persons age 15 or older. According to the MOI, police investigated 529 cases of violence against women and children, including child sexual exploitation, rape, debauchery, and human trafficking. The investigations resulted in the arrest of 582 offenders, of whom 46 were arrested for cross border and domestic trafficking. However, NGOs continued to report the general failure of law enforcement and other government officials to act on tip‑offs.
The Ministries of Interior, Women's Affairs, and Justice had primary responsibility to combat trafficking in persons. In April the government established a National Task Force to serve as an interministerial antitrafficking coordination body. The task force included an oversight body involving the top government officials. There was a Department of Anti‑Human Trafficking and Juvenile Protection, and the MOI operated specialized antitrafficking divisions in all provinces and municipalities. While the government arrested and prosecuted traffickers and continued its support for prevention and protection programs through collaboration with foreign and domestic NGOs and international organizations, its antitrafficking efforts continued to be hampered by corruption and a weak judicial system. It was widely believed that some law enforcement and other government officials received bribes that facilitated the sex trade and trafficking in persons.
On March 16, the Sihanoukville Municipal Court acquitted the owner of O Pi Guesthouse and an employee of all charges and convicted another employee of a lesser charge of pimping, sentencing her to a two‑year imprisonment plus a three‑year suspended sentence. Erratic official behavior during the trial and the light penalty raised concerns that there was an exchange of bribes in return for light treatment of the case.
On March 21, the Phnom Penh Municipal Court acquitted Meng Say, former chief of the Phnom Penh antitrafficking unit, who was suspended in 2006 for extorting money from South Korean nationals. One police officer remained in jail in connection with the 2006 Phnom Penh Municipal Court case of three police officers sentenced to five to seven years in prison for trafficking‑related corruption. Two of the officers appealed their cases; there was no court action on the appeal, and the two officers had not started serving their sentences.
On August 1, the Sihanoukville Municipal Court tried a pimping case but acquitted the suspect, citing lack of evidence. Observers reported irregularities in the case hearing, and the case was under appeal at year's end.
On August 9, a royal decree directed the Supreme Council of the Magistracy to dismiss Appeals Court President Ly Vouch Leng. The directive was issued for her alleged acceptance of bribes in exchange for the release of human traffickers who were running the Chhay Hour II brothel in Phnom Penh. Three Supreme Council of Magistracy officials were also removed in connection with the case, and three appeals court judges and one deputy prosecutor received official letters of reprimand. Ly Vouch Leng was transferred to an unknown position in the Ministry of Justice; no charges were brought against her. The Ministries of Interior and Justice reported their investigations continued.
A legal advocacy NGO reported that five trafficking cases went to trial during the year, resulting in two convictions. For the same period, the MOI reported five convictions on human trafficking charges with sentences ranging from four to 15 years in prison. The Phnom Penh Municipal Court reported 49 convictions of human trafficking offenders from January to October. Police, court officials, and judges often did not separate victims from perpetrators during raids, arrests, and trials. In some cases officials spoke and acted as though victims were perpetrators. During a March 9 Sihanoukville trafficking trial, the presiding judge spoke harshly to underage trafficking victims in the courtroom and acquitted two of the alleged perpetrators.
The Ministry of Social Affairs and Youth Rehabilitations (MOSAVY) referred trafficking victims to NGOs, which provided most assistance to victims. The government participated as a partner in a number of these efforts; however, its contributions were severely hampered by limited resources. NGOs provided intake screening services to identify trafficking victims. Some victims were encouraged by NGOs and the MOI to file complaints against perpetrators; however, in the general climate of impunity, victim protection was problematic, and victims were known to be intimidated into abandoning their cases. Social stigma against women who were prostitutes, victims of sexual assault, or victims of sex trafficking made it difficult for victims to reintegrate into families, communities, and society.
The trafficking law contains no provisions to protect foreign victims from being charged under immigration laws, but during the year there were no reported cases of trafficking victims being treated as illegal immigrants. The MOSAVY worked with the IOM to repatriate trafficked victims from Thailand and Vietnam to Cambodia, and from Cambodia to Vietnam. However, repatriation to Vietnam continued to be a long and arduous process.
The MOSAVY repatriated from Thailand, Vietnam, and Malaysia 845 child and adult victims, as well as others vulnerable to becoming victims, and reintegrated them with their families. With financial and technical support from the IOM, the MOSAVY repatriated eight trafficked Vietnamese girls to Vietnam.
Both the government and international donors had programs to prevent child labor or remove children from labor. The country is a signatory to the Coordinated Mekong Ministerial Initiative against Trafficking, whose activities include ensuring the legal, social, and community protection of victims of trafficking; strengthening law enforcement capacity to combat trafficking; and building a comprehensive response involving all relevant ministries. Several ministries--including the Ministry of Women's Affairs and the Ministry of Tourism--had antitrafficking initiatives to reduce child labor. Donors supported programs to combat child labor implemented by the ILO and World Education, among others.
The MOSAVY worked with the UN Children's Fund (UNICEF) and local NGOs to manage community‑based networks aimed at preventing trafficking.
Persons with Disabilities
There is no law explicitly prohibiting discrimination against persons with disabilities. The government does not require that buildings or government services be accessible to persons with disabilities. The government prohibits persons with disabilities from being teachers in public schools. On October 1, the government signed the UN Convention on the Rights of Persons with Disabilities.
Programs administered by various NGOs brought about substantial improvements in the treatment and rehabilitation of persons who had lost limbs, but they faced considerable societal discrimination, especially in obtaining skilled employment.
There are no legal limitations on the rights of persons with disabilities to vote or participate in civic affairs, but the government did not make any concerted effort to assist them in becoming more civically engaged. The MOSAVY is responsible for making policy to protect the rights of persons with disabilities and for rehabilitation and vocational skill training for persons with disabilities.
The rights of minorities under the 1996 nationality law are not explicit; constitutional protections are extended only to "Khmer people." Citizens of Chinese and Vietnamese ethnicity constituted the largest ethnic minorities. Ethnic Chinese citizens were accepted in society, but animosity continued toward ethnic Vietnamese, who were seen as a threat to the nation and culture. Some groups continued to make strong anti‑Vietnamese statements. They complained of political control of the CPP by the Vietnamese government, border encroachment, and other problems for which they held ethnic Vietnamese at least partially responsible.
The government often ignored efforts by indigenous communities to protect their ancestral lands and natural resources. In spite of the 2001 land law, which calls for the registration of communal lands of indigenous people, little was done to implement communal land titling. NGOs called for a moratorium on land sales and land concessions affecting indigenous communities. International and local NGOs were active in educating the indigenous communities about their land rights and providing legal representation in disputes.
On March 9, more than 200 indigenous villagers in Stung Treng Province protested the clearing of community forest land by four companies to which the government allegedly illegitimately granted timber concessions. The land had long been used by indigenous villagers for subsistence farming, hunting, and resin production. In May provincial authorities created a special committee to resolve the problem, but the committee did not take any action.
On March 15, more than 500 Jarai indigenous families in Ratanakiri Province demanded the removal of local officials who they alleged were involved in the fraudulent sale of more than 3,000 acres of their communal land. In September 2006 the villagers learned their land had been sold when they saw workers demarcating it as private property. The villagers submitted a complaint to provincial authorities, but authorities did not respond to the complaint.
In early August a Ratanakiri provincial official prevented Tampoun indigenous villagers from burying their dead on land that had served as their traditional burial ground since 1979. The provincial court stated it would arrest anyone who tried to bury bodies there, claiming the land belonged to the provincial finance department director. Authorities sanctioned a new burial ground approximately 500 yards from the traditional plot. The villagers enlisted the help of an NGO and planned to file a suit with the provincial court.
During the August 9 commemoration of the International Day for Indigenous Peoples, the UN High Commissioner for Human Rights noted the government's failure to protect and implement the rights of indigenous people to their lands, territories, and natural resources. The UN commissioner called for swift action to halt land grabbing in tribal areas, particularly the growing number of economic land concessions and mining licenses granted without community consultation.
Other Societal Abuses and Discrimination
Societal discrimination against those infected with HIV/AIDS remained a problem in rural areas; however, discrimination was moderated by HIV/AIDS awareness programs. There was no official discrimination against those infected with HIV/AIDS. There were no reported cases of sexual orientation discrimination in employment, housing, statelessness, or access to education or health care. However, homosexuality was typically treated with fear and suspicion, and there were few support groups based on sexual orientation where such cases could have been reported.
Section 6 Worker Rights
a. The Right of Association
The labor law provides private-sector workers in the formal economy the right to join the trade union of their choice without prior authorization. However, the government's enforcement of this right was selective. Membership in trade unions or employee associations is not compulsory, and workers are free to withdraw from such organizations, although a few unions attempted to intimidate workers who wanted to withdraw. Unions may affiliate freely, but the law does not address explicitly their right to affiliate internationally. While the law applies to foreign workers, it does not apply to civil servants, including teachers, judges, and military personnel, or to workers in the informal sector. Personnel in the air and maritime transportation industries are not entitled to the full protections of the law but are free to form unions.
The vast majority of the country worked in the informal sector, primarily as subsistence rice farmers, vendors, or skilled or unskilled laborers. Only a small fraction, estimated at less than 1 percent, of the labor force was unionized. Unions were concentrated in the garment and footwear industries, where approximately 40 to 50 percent of the 350,000 workers were union members. The Cambodian Tourism and Service Workers Federation represented 4,000 hotel, casino, and airport workers. Of the 31 national labor federations and confederations, 26 were allied with the government, four were independent, and one had pro‑opposition leanings.
The Cambodia Independent Teachers Association (CITA), registered as an "association" due to prohibitions on public sector unions, represented 8,150 of the country's 89,000 teachers. CITA marches and other protests were often forbidden, although the union reported no direct government interference in day‑to‑day activities. Some members feared that CITA affiliation could hamper their chance of career advancement, according to union officials. Another public sector association, the Cambodian Independent Civil Servants' Association (CICA), represented approximately 500 officials from ministries, provincial departments, and commune councils, out of approximately 100,000 civil servants nationwide. CICA leaders alleged that fears of harassment or demotion prevented other civil servants from joining.
Some independent and pro‑opposition unions and federations complained of unnecessary delays and costs in registering with the government.
Unions were generally seen as slowly gaining strength, but many were not able to adequately represent member interests due to insufficient resources, training, and experience. In addition, corruption plagued unions, employers, and government officials, hampering legitimate industrial relations. Violence, harassment, and intimidation between rival unions were common.
On February 24, two unidentified men shot and killed local union leader Hy Vuthy as he left the Suntex garment factory after completing his night shift. Since 2005 Suntex and Bright Sky factories, which share a compound, have been the scene of fierce interunion rivalry and violence. Hy Vuthy was a member of the country's largest union, the Free Trade Union of Workers in the Kingdom of Cambodia (FTUWKC), which alleged that he might have been killed because of his labor work. No suspects were arrested. Two other FTUWKC leaders--national FTUWKC president Chea Vichea and local union leader Ros Sovannareth--were killed in 2004.
In some factories persons employed in management appeared to have established their own unions, supported promanagement unions, or compromised union leaders. Union leaders from across the political spectrum complained that the Khmer Youth Federation of Trade Unions habitually threatened strikes to extort money from management and threatened and harassed workers from other unions. Independent union leaders complained that the progovernment Cambodian Coalition of Trade Unions frequently intervened in the affairs of other unions, extorted money from management in exchange for discouraging workers from conducting legal strikes and demonstrations, and threatened rival union leaders.
Enforcement of the right of association and freedom from antiunion discrimination was poor. Government enforcement was hampered by a lack of political will and by confused financial and political relationships with employers and union leaders. The government also suffered from a lack of resources, including trained, experienced labor inspectors, in part because it did not pay staff adequate salaries. The Ministry of Labor and Vocational Training (MOLVT) often decided in favor of employees but rarely used its legal authority to penalize employers who defied its orders, instead referring many cases to an arbitration council.
There were credible reports of antiunion harassment by employers, including the dismissal of union leaders, in garment factories and other enterprises. Employers sometimes used the courts to dismiss or punish union leaders. In two cases union leaders were charged with inciting workers to strike, destroying private property, and attempting to incite workers to commit assault. At year's end the cases were pending. On several occasions dismissed union leaders accepted cash settlements after unsuccessfully appealing to the government to enforce laws requiring their reinstatement. At other times the government upheld labor rights. For example, the MOLVT formally warned 1,032 companies of legal violations, fined 10 companies, and charged five companies with violation of the labor law and regulations. The MOLVT sent 83 cases of unresolved labor disputes to the Arbitration Council.
b. The Right to Organize and Bargain Collectively
The law provides for the right to organize and bargain collectively, but the government's enforcement of these rights was inconsistent. Wages were generally set by market forces, except in the case of civil servants, whose wages were set by the government. Garment-sector workers were guaranteed a minimum wage of $50 per month (210,000 riels).
During the year there were 19 collective bargaining agreements registered with the MOLVT. Most were conciliation agreements that did not meet international collective bargaining standards. Only six genuine collective bargaining agreements existed within the garment industry, 10 at hotels, and one covering contract workers at the two international airports.
A 2001 regulation establishes procedures to allow unions to demonstrate that they represent workers for purposes of collective bargaining. The regulation also establishes requirements for employers and unions regarding collective bargaining and provides union leaders with additional protection from dismissal. The Bureau of Labor Relations is responsible for facilitating the process of union registration and certification of "most representative status" for unions, which entitles a union representing an absolute majority of workers in a given enterprise to represent all of the workers in that establishment. However, the "most representative" registration process was considered cumbersome, and international observers reported that government lists of "most representative unions" included management‑controlled unions and unions whose "most representative" status should have expired years before. The government began reexamining its "most representative" certification process with support from international organizations and a diplomatic mission.
The law provides for the right to strike and protects strikers from reprisal. The law stipulates that strikes can be held only after several requirements have been met, including the failure of other methods of dispute resolution (such as negotiation, conciliation, and arbitration), a secret-ballot vote of union membership, and a seven‑day advance notice to the employer and the MOLVT.
The MOLVT reported that 82 strikes occurred during the year. International observers, employers, and many union leaders agreed that almost no strikes fulfilled all prestrike legal requirements. Other unions complained that a severe lack of MOLVT involvement led to industrial strikes.
The government allowed most strikes held at factories but denied worker requests to hold protest marches outside of the factory district. Police intervention in strikes generally was minimal and restrained, even in those cases where property damage occurred. Police presence at the few marches that occurred tended to be excessive and often included a specialized police intervention unit.
On May 4, four workers and five security guards were injured when security guards inside the L.A. garment factory tried to prevent workers from leaving the factory to strike. The strike, which involved 4,000 workers, began two days earlier when workers demanded that a manager be dismissed for insulting and mistreating them.
On May 21, more than 100 provincial police officers violently dispersed approximately 200 striking workers at the River Rich garment factory in Kandal Province. Workers said police beat several protesters, but police denied causing any injuries. The strike began when management reneged on a promise to rehire 10 union activists whose contracts had not been renewed. Following the strike, management threatened to sue three union representatives for inciting workers to hold an illegal strike and for discrimination but later dropped the case.
On November 29, approximately 200 police officers violently dispersed a strike by more than 2,000 Fortune garment factory workers who protested reductions in bonuses and short‑term contracts. Police accused the strikers of blocking a road, creating disorder, throwing stones, and injuring police officers. Workers said their strike was peaceful and that violent police repression--including firing guns into the air and using tear gas--resulted in three injured workers. Police detained four workers, releasing them later the same day.
There were no developments in the August 2006 case of three factory‑level union leaders affiliated with the FTUWKC convicted of charges of illegal human confinement. After spending one month and four days in jail in 2006, the workers were released.
In spite of legal provisions protecting strikers from reprisals, there were credible reports that workers were dismissed on spurious grounds after organizing or participating in strikes. While most strikes were illegal, participating in an illegal strike was not by itself a legally acceptable reason for dismissal. In some cases strikers were pressured by employers to accept compensation and leave their employment. There are potential remedies for such dismissals, although none was particularly effective. The MOLVT can issue reinstatement orders, but these often provoked management efforts to pressure workers into resigning in exchange for a settlement. Collective disputes, such as when multiple employees are dismissed, can be brought before the Arbitration Council for a nonbinding decision. Individual disputes can be bought before the courts, although the judicial system was neither impartial nor transparent. Some unions urged the government to expand the role of the Arbitration Council to include individual and collective interest disputes and to make its decisions binding.
There continued to be confusion about the overlapping roles of labor unions and elected shop stewards. According to regulation, trade union leaders have roles comparable to those of shop stewards, and certain union officers have protection from dismissal within an enterprise. However, employers did not always respect these protections.
There were no special laws or exemptions from regular labor laws in export processing zones (known as special economic zones).
c. Prohibition of Forced or Compulsory Labor
The law prohibits forced or compulsory labor, including by children, but there were reports that such practices occurred, almost exclusively in the informal sector. There were reports of isolated cases of forced labor by domestic servants. Forced child labor was a serious problem in the commercial sex industry.
Involuntary overtime remained widespread. Under the law, legal overtime work cannot exceed two hours daily and must be voluntary; however, in practice overtime was often extended beyond the legal limit, and employers used coercion to force employees to work. Workers often faced fines, dismissal, or loss of premium pay if they refused to work overtime.
d. Prohibition of Child Labor and Minimum Age for Employment
The government has adopted laws to protect children from exploitation in the workplace; however, enforcement was often weak. The law establishes 15 years as the minimum age for employment and 18 years as the minimum age for hazardous work. The law permits children between 12 and 15 to engage in "light work" that is not hazardous to their health and does not affect school attendance. A 2006 study by the World Bank, the ILO, and UNICEF estimated that there were 1.5 million children engaged in illegal labor, including 750,000 children younger than 12 years, 500,000 children ages 12 to 14 engaged in "nonlight" economic activity, and more than 250,000 children ages 15 to 17 working in prohibited hazardous sectors or working more than 43 hours per week.
No aspect of the law prohibiting child labor was adequately enforced in the formal employment sector. No employer was prosecuted for violating laws against child labor. The MOLVT has responsibility for child labor issues in both the formal and informal sectors of the economy, but its labor inspectors played no role in the informal sector or in enforcing the law in illegal industries. Within the formal sector, labor inspectors conducted routine inspections of some industries, such as garment manufacturing (where the incidence of child labor is negligible), but in some industries with the highest child labor risk, labor inspections were entirely complaint driven.
The constitution prohibits forced or bonded child labor; however, forced child labor was a serious problem in the commercial sex industry. Law enforcement agencies failed to combat child prostitution in a sustained, consistent manner. Widespread corruption, lack of transparency, inadequate resources, and staffing shortages remained the most challenging obstacles.
e. Acceptable Conditions of Work
The law requires the MOLVT to establish a garment‑sector minimum wage based on recommendations from the Labor Advisory Committee. There was no minimum wage for any other industry. The minimum wage for the sector was $45 to $50 (189,000 to 210,000 riels) per month. Garment-sector employers almost universally paid regular workers at least the minimum wage, although casual workers were often paid less. Garment workers earned an average wage of $70 to $80 (294,000 to 336,000 riels) per month, including overtime and bonuses. Prevailing monthly wages in the garment sector and many other professions were insufficient to provide a worker and family with a decent standard of living, although garment-sector wages were generally higher than wages in the informal economy. Civil service salaries also were insufficient to provide a decent standard of living, requiring government officials to secure outside sources of income, in many cases by obtaining second jobs or accepting bribes.
On June 8, despite strong protests from some unions, the National Assembly amended the labor law to establish a nightshift rate of 130 percent of daytime wages. Before this amendment, customary practice was to pay nightshift workers 200 percent of daytime wages, and few factories operated night shifts due to the high salary cost.
The law provides for a standard legal workweek of 48 hours, not to exceed eight hours per day. The law stipulates time‑and‑a‑half for overtime and double time if overtime occurs at night, on Sunday, or on a holiday. Employees are allowed to work up to two hours of overtime each day. However, the government did not enforce these standards effectively. Workers reported that overtime was frequently excessive and sometimes mandatory. Similarly, outside the garment industry, regulations on working hours were rarely enforced.
The law states that the workplace should have health and safety standards adequate to ensure workers' well‑being. The government enforced existing standards selectively, in part because it lacked trained staff and equipment. Work‑related injuries and health problems were common. Most large garment factories producing for markets in developed countries met relatively high health and safety standards as conditions of their contracts with buyers. Working conditions in some small‑scale factories and cottage industries were poor and often did not meet international standards. Penalties are specified in the law, but there are no specific provisions to protect workers who complain about unsafe or unhealthy conditions. Workers who removed themselves from unsafe working conditions risked loss of employment. |
Name ID 1199
The first professional hunters came in 1913. They found the wildlife plentiful, especially the lions, but saw no elephants. Seven years later, an American arrived in a strange new contraption known as a Ford motor-car and news of the wonders of the Serengeti had reached the outside world. Because the Hunting of lions made them so scarse (they were considered 'vermin'), it was decided to make a partial Game Reserve in the area in 1921 and a full one in 1929. With the growing awareness of the need for conservation, it was expanded and upgraded to a National Park in 1951. Eight years later the Ngorongoro Conservation Area was established in the south-east as a separate unit.
Arusha: A Brochure of the Northern Province and its Capital Town
Page Number: 13-15-17
Extract Date: 1929
It is safe to say that Tanganyika holds a front place among our East African Colonies for the number and variety of its game animals. The belt from Tanga through to Lake Victoria is where game is most numerous. There is an abundance of the commoner antelope, and in certain parts the rarer species such as the Greater and Lesser Kudu, Gerenuk, etc., are still fairly plentiful. Big game like the Elephant, Rhinocerous, Lion and Buffalo, all of which hold for the hunter a new thrill and experience, are to be found in this area in such a variety of country and cover that the Hunting of no two animals is ever alike.
Here the hunter passes through most interesting country; Kilimanjaro with its snow-capped dome, running streams and dense forests, across the plains to the Natron Lakes and the Great Rift Wall with its volcanic formation and on to the great Crater, Ngorongoro. In his travels he will come into contact with some of the most interesting and picturesque tribes that inhabit Africa such as the Masai, Wambulu, etc., each with their own quaint customs and histories.
The Ngorongoro Crater, the greatest crater in the world, measuring approximately 12 miles in diameter, seen from the Mbulu side, is a delight to the eye with its teeming herds of game ; Wildebeest alone running into tens of thousands. This scene conveys to one the idea of a great National Park. Nature has provided the crater with a precipitous rock fence for tns most part and with lakes and streams to slake the thirst of the great game herds which inhabit it. The unalienated part of the crater is now a complete game reserve in which a great variety of game is to be found such as Rhinoceros, Hippopotamus, Lion, and all the smaller fry. The Elephant although not in the crater is to be found in the forests nearby.
The Serengetti Plains lying away to the northwest of the crater holds its full share of animal life and here the sportsman has the widest possible choice of trophies. The Lion in this area holds full sway and is still to be seen in troops of from ten to twenty. Recently, Serengetti and Lion pictures have become synonymous. The commoner species of game are here in abundance and the plains are second only to the crater for game concentration. The country lying between the Grumeti River - Orangi River and the Mbalangeti from Lake Victoria to the Mou-Kilimafetha Road has recently been declared a game reserve.
Game animals that inhabit the northern area are well protected an'd their existence is assured to posterity by the great game sanctuaries and regulations which govern the Hunting or photographing of game.
In the Northern area there are six complete reserves and two closed areas. These are as follows:
(2) Mount Meru.
(3) Lake Natron
(4) Northern Railway.
The closed areas are :
Pienaar's Heights, near Babati and Sangessa Steppe in the Kondoa district. The boundaries for these are laid down in the Game Preservation Ordinance No. 41 of 1921. There are, however, vast areas open to the hunter and the abovementioned sanctuaries do not in any way detract from the available sport which the Northern Tanganyika has to offer.
The following game licences are now in force (Shillings)
:Visitor's Full Licence - 1500
Visitor's Temporary Licence (14 days) - 200
Resident's Full Licence - 300
Resident's Temporary Licence (14 days) - 60
Resident's Minor Licence - 80
Giraffe Licence - 150
Elephant Licence 1st. - 400
2nd. - 600
To hunt the Black Rhinoceros in the Northern Province it is now necessary to hold a Governor's Licence, the fee for which is 150/-. This entitles the holder to hunt one male Rhinoceros. Elephant, Giraffe, and Rhinoceros Licences may only be issued to holders of full game licences.
Now that the Railway is through to Arusha it is not too much to hope that with the assistance of a healthy public opinion the Sanya Plains may become restocked with game which would be a great source of interest and an attraction to the traveller visiting these parts.
Extract Date: 1931
Herne, Brian White Hunters: The golden age of African Safaris
Page Number: 375
Extract Date: 1965
Safari Hunting in East Africa was forever changed by the masterly blueprint of Brian Nicholson, a former white hunter turned game warden. The disciple and successor of C.I. P. Ionides, the "Father of the Selous game reserve," Nicholson conceived a plan for administering Tanzania's expansive wildlife regions. In 1965 he changed most of the vast former controlled Hunting areas, or CHAs, into Hunting concessions that could be leased by outfitters from the government for two or more years at a time. Nicholson also demarcated the Selous game reserve's 20,000 square miles of uninhabited country into 47 separate concessions. Concessions were given a limited quota of each game species, and outfitters were expected to utilize quotas as fully as possible, but not exceed them.
Nicholson's plan gave outfitters exclusive rights over Hunting lands, providing powerful incentives for concession holders to police their areas, develop tracks, airfields, and camps, and, most importantly, preserve the wild game. When the system was put into effect, it was the larger outfitting organizations - safari outfitters who could muster the resources to bid and who had a clientele sufficient to fulfill the trophy quotas Nicholson had set (done in order to provide government revenue by way of fees for anti-poaching operations, development, and research) - that moved quickly to buy up the leases on the most desirable blocks of land. Smaller safari companies who could not compete on their own banded together and formed alliances so that they, too, could obtain Hunting territories.
Herne, Brian White Hunters: The golden age of African Safaris
Page Number: 389
Extract Date: 1973 Sep 7
By the end of 1973 Kenya was the sole remaining tourist destination in East Africa. While the neighboring country of Uganda was still in the throes of military anarchy, Tanzania surprised the world on September 7 by issuing an overnight ban on all Hunting and photographic safaris within its territory. Government authorities moved quickly to seize and impound foreign-registered Land Cruisers, supply trucks, minibuses, aircraft, and equipment.
The stunned collection of safari clients as well as sundry mountain climbers, bird-watchers, and beachcombers who had been visiting the country at the time of the inexplicable edict were summarily escorted to Kilimanjaro airport outside of Arusha to await deportation. The residue of tourists stranded without flights were trucked to the northern town of Namanga where they were left on the dusty roadside to cross into Kenya on foot. All tourist businesses, including the government-owned Tanzania Wildlife Safaris, were closed down. No government refunds were ever made to tourists or to foreign or local safari outfitters
Africa News Online
Extract Date: 2000 June 5
Panafrican News Agency
Frequent acrimony, currently depicting the relationship between game Hunting companies and rural communities in Tanzania, will be a thing of the past after the government adopts a new wildlife policy.
Designated as wildlife management areas, the communities will benefit from the spoils of game Hunting, presently paid to local authorities by companies operating in those areas.
The proposed policy seeks to amend Tanzania's obsolete Wildlife Act of 1975, and, according to the natural resources and tourism minister, Zakia Meghji, 'it is of utmost priority and should be tabled before parliament for debate soon'.
She said the government would repossess all Hunting blocks allocated to professional hunters and hand them over to respective local authorities.
In turn local governments, together with the communities, would be empowered to allocate the Hunting blocks to whichever company they prefer to do business with.
'Guidelines of the policy are ready and are just being fine-tuned,' she said.
Communities set to benefit from this policy are chiefly those bordering rich game controlled areas and parks. They include the Maasai, Ndorobo, Hadzabe, Bahi, Sianzu and Kimbu in northeastern Tanzania.
Members of these communities are often arrested by game wardens and fined for trespassing on game conservation areas. As a result, they have been extremely bitter about being denied access to wildlife resources, which they believe, naturally, belong to them.
Under the new policy, Meghji said, the government will ensure that people undertake increased wildlife management responsibilities and get benefits to motivate them in the conservation of wildlife resources.
The East African
Extract Author: John Mbaria
Extract Date: February 4, 2002
KENYA COULD end up losing 80 per cent of its wildlife species in protected areas bordering Tanzania to hunters licensed by the Tanzanian government.
The hunters have been operating for about a decade in a section of the migratory route south from Kenya to Tanzania's Serengeti National Park.
They shoot large numbers of animals as they move into the park during the big zebra and wildebeest migration between July and December.
There are fears that the Maasai Mara National Park and most of Kenya's wildlife areas bordering Tanzania could lose much of their wildlife population, threatening the country's Ksh20 billion ($256 million) a year tourism industry.
Kenya banned Hunting in 1977 but the sport is legal in Tanzania, where it is sold as "Safari Hunting."
"The product sold is really the experience of tracking and killing animals, the services that go with this and the prestige of taking home the trophies," says a policy document from the Tanzania Wildlife Corporation (Tawico).
Tanzania wildlife officials said wild animals that cross over from Kenya are hunted along their migratory routes in the Loliondo Game Controlled Area in Ngorongoro district of Arusha region, 400 km northwest of Arusha. The area was designated by the British colonial power as a sports Hunting region for European royalty.
The officials said the area is now utilised by a top defence official from the United Arab Emirates (UAE), trading as Ortelo Business Company (OBC), through a licence issued in 1992 by former Tanzania President Ali Hassan Mwinyi. The permit allows the company to hunt wild game and trap and take some live animals back to the UAE.
Safari Hunting earns big money for the Tanzanian government, which charges each hunter $1,600 a day to enter the controlled area.
A hunter is also required to pay fees for each kill, with an elephant costing $4,000, a lion and a leopard $2,000 each and a buffalo $600. The document has no quotation for rhinos.
The sport is organised in expeditions lasting between one and three weeks in the five Hunting blocks of Lake Natron Game Controlled Area, Rungwa Game Reserve, Selous Mai, Selous U3 and Selous LU4.
For the period the hunters stay in each of the Hunting blocks, they pay between $7,270 and $13,170 each. Part of this money is shared out among the many Ujamaa villages, the local district councils and the central government.
Although Tawico restricts the number of animals to be culled by species, poor monitoring of the activities has meant indiscriminate killing of game.
"Some of the animals are snared and either exported alive or as meat and skins to the United Arab Emirates and other destinations," local community members told The EastAfrican during a recent trip to the area.
They claimed the hunters were provided with "blank Hunting permits," giving them discretion over the number of animals to be hunted down. Kenya wildlife conservation bodies are concerned that big game Hunting in the Ngorongoro area is depleting the wildlife that crosses the border from Kenya.
"Kenya is losing much of its wildlife to hunters licensed by the Tanzanian government," the chairman of the Maasai Environmental and Resource Coalition (MERC), Mr Andrew ole Nainguran, said. MERC was set up in 1999 to sensitise members of the Maasai community in Kenya and Tanzania to the benefits of wildlife conservation.
Kenya and Tanzania wildlife authorities have regularly discussed the problem of security and poaching in Arusha. However, the KWS acting director, Mr Joe Kioko, said legalised Hunting has never been discussed in any of the meetings.
The hunters are said to fly directly from the UAE to the area using huge cargo and passenger planes which land on an all-weather airstrip inside the OBC camp. The planes are loaded with sophisticated Hunting equipment, including four-wheel drive vehicles, weapons and communication gadgets.
On their way back, the planes carry a variety of live animals, game trophies and meat. Employees at the camp said the hunters are sometimes accompanied by young Pakistani and Filipino women.
The International Fund for Animal Welfare regional director, Mr Michael Wamithi, said Kenya and Tanzania should discuss the negative impact of the sport Hunting on Kenya's conservation efforts.
"The two neighbours have a Cross-Border Law Enforcement Memorandum of Understanding where such issues could be dealt with."
Kenya seems to be alone in adhering to strict protection of wildlife, a policy famously demonstrated by President Daniel arap Moi's torching of ivory worth $760,000 in 1989.
Although the country has made significant progress in securing parks from poachers, it is yet to embrace a policy on "consumptive utilisation" of animals advocated by Kenyan game ranchers and Zimbabwe, which wants the international trade ban on ivory lifted.
The animals in the Hunting block have been reduced to such an extent that the OBC camp management has been spreading salt and pumping water at strategic places to attract animals from Serengeti and the outlying areas.
"We will not have any animals left in the vicinity unless the Hunting is checked," a local community leader, Mr Oloomo Samantai ole Nairoti, said, arguing that the area's tourism economy was being jeopardised.
Mysterious fires in the area to the south of Serengeti have also forced animals to seek refuge in the Hunting blocks.
Locals said the camp is exclusively patronised by Arab visitors. The camp is usually under tight security by Tanzanian police.
The permit granted by Mr Mwinyi has raised controversy in Tanzania and was at one stage the subject of a parliamentary probe committee because members of UAE's royal family were not entitled to the Hunting rights in the country.
"Only presidents or monarchs are entitled to hunt in the area," an official said, adding that the UAE royal family had abused their permit by killing animals outside their given quotas or specified species.
The government revoked the licence in 1999 after realising that OBC was airlifting many wild animals to the Middle East, only to renew the permit in 2000. The current permit runs until 2005.
The withdrawal of the permit followed the recommendations of a 1994 parliamentary probe commission set up to "investigate the Hunting behaviour" of the UAE company.
Sources said permanent Hunting is prohibited in the Loliondo Game Controlled Area for fear of depleting animals from the four parks, which host the bulk of the region's tourist resorts.
The area is in a natural corridor where wild animals cross while roaming between the Ngorongoro Conservation Area and Serengeti National Park in Tanzania and Maasai Mara Game Reserve and Amboseli National Park in Kenya.
The late founding president of Tanzania, Mwalimu Julius Nyerere, took to himself the powers to issue Hunting permits for Loliondo when Tanzania became independent in 1961, but he never granted any.
After obtaining the permit, the UAE hunters created Hunting blocks in the area covering over 4,000 sq km.
No other Hunting companies have been granted permits, the source said.
The UAE royal family has donated passenger aircraft to the Tanzania army and a number of vehicles to the Wildlife Division.
The 1974 Wildlife Act set up five categories of wildlife conservation areas.
These are national parks, game reserves, partial game reserve, open areas and Ngorongoro Conservation Area. Hunting is prohibited in the national parks and Ngorongoro Conservation Area, but allowed in other areas during the seasonal Hunting period from July to December.
Additional reporting by Apolinari Tairo in Dar es Salaam
Tomlinson, Chris Big game hunting threatening Africa
Extract Author: Chris Tomlinson
Extract Date: 2002 03 20
Loliondo GAME CONTROL AREA, Tanzania - At a dirt airstrip in rural Tanzania, a desert camouflaged cargo plane from the United Arab Emirates air force taxis up to pallets stacked with large coolers full of game meat, the harvest of a successful Hunting season.
As Tanzanian immigration and customs officials fill out documents under a thatched shelter, brand-new, four-wheel-drive trucks and dune buggies drive to and from a nearby luxury campsite, the base for one of Tanzania's most expensive - and secretive - game Hunting operations, Otterlo Business Corp.
Hundreds of members of Arab royalty and high-flying businessmen spend weeks in the Loliondo Game Control Area each year Hunting antelope, lion, leopard and other wild animals. The area is leased under the Otterlo name by a member of an emirate royal family who is a senior officer in the UAE defense ministry.
While neighboring Kenya outlawed big game Hunting in 1978, the Tanzanian government says Hunting is the best use of the land and wildlife. But villagers and herders say big money has led government officials to break all the Hunting rules, resulting in the destruction of most of the area's non-migratory animals and putting East Africa's most famous national parks under threat.
Loliondo is on the main migratory route for wildlife north of Ngorongoro Crater, east of Serengeti National Park and south of Kenya's Masai Mara National Reserve. The summer Hunting season coincides with the migration of wildebeest and zebra through the area, where they eventually cross into the Serengeti and the Masai Mara. Predatory animals follow the migration.
During the colonial era, Loliondo was set aside for European royalty as a Hunting area. Since independence, Loliondo has remained a Hunting reserve, but it is supposed to be managed by area residents for their benefit.
Local leaders, who refuse to speak publicly because they fear retribution, say they have not been consulted about the lease that was granted in 1995 by national officials in Tanzania's political capital, Dodoma. They say government officials have tried to silence criticism.
"The lease was given by the government and the Maasai landowners were not involved," said one Maasai leader. "All the resident animals have been killed ... (now) they carry out Hunting raids in the Serengeti National Park, but the government closes its eyes."
Maasai warriors told The Associated Press that hunters give cash to anyone who can lead them to big game, especially leopards. They also said that Otterlo officials have begun pumping water into some areas to attract more animals and that what the warriors call suspicious fires in the Serengeti have caused animals to move into Loliondo.
An Otterlo manager, who gave his name only as Khamis, initially agreed to an interview with AP but later did not return repeated phone calls.
In an interview with the newspaper, The East African, Otterlo managing director Juma Akida Zodikheri said his company adheres to Tanzanian law, and he denied hunters killed animals indiscriminately. He said the owner of the company is Maj. Gen. Mohammed Abdulrahim al Ali, deputy defense minister of the UAE.
While Tanzania has strict rules on game Hunting, Maasai who have worked at the lodge say guests are never told of the limits and hunt as much as they want. Tanzanian officials deny that.
Col. A.G.N. Msangi, district commissioner for Ngorongoro District, said all applicable rules are enforced. He accused the Maasai of rumor-mongering in an effort to discredit Otterlo.
The company "is following the system the government wants," Msangi said. "OBC has invested more money here than any other company in the district."
Msangi said Hunting companies request permission to kill a certain number of animals. Once the request is approved by wildlife experts at the Ministry of the Environment, the company pays a fee based on that number whether they actually kill the animals or not, he said.
"We have police and ministry people making sure they don't exceed what they have paid for," Msangi said. The tourists are also required to employ professional hunters to ensure no female or young animals are killed, he added.
Compared to the numbers in Serengeti National Park, very few large animals were seen during a three-hour drive through Loliondo. But without any independent survey of the animal population, it is impossible to know whether Msangi's conservation efforts are working.
Msangi described his main duty as balancing the needs of people, animals and conservation. He said not only does Hunting revenue finance wildlife conservation, but Otterlo, like most tourism companies, also makes charitable donations to help pay for schools and development projects and it provides badly need jobs.
Also appeared in http://www.washtimes.com/world/20020801-22110374.htm
1 Aug 2002
Internet Web Pages
Extract Author: Lifer
Extract Date: April 16 2002
Posted - April 16 2002 : 20:53:22
The East African Newspaper of 4-10 February 2002 carried an article titled "Game Carnage in Tanzania Alarms Kenya", written by John Mbaria with supplement information from Apolinari Tairo of Dar es Salaam. The article was on The Ortello Business Hunting Company, which started to hunt in the Loliondo Game Controlled Area in 1992.
The following are issues raised in the article:
a) Hunting activities carried out in Liliondo Game Controlled Area near the Tanzania / Kenyan border causes loses of 80% of the Kenyan wildlife.
b) Hunting is conducted in the migratory route in the south between Kenya and Serengeti National Park. The animals are hunted during the migratory period as they move to Kenya and on their way back to Tanzania in July to December.
c) Hunting is threatening the Kenyan tourism industry, which earns the country USD 256.0 annually.
d) The Hunting kills animals haphazardly, without proper guidance and monitoring of actual number of animals killed and exported outside the country.
e) Airplanes belonging to Ortello Business Corporation (OBC) carry unspecified type of live animals and birds from Loliondo on their way back to UAE. Further more, the air planes fly directly in and out of Loliondo without stopping at Kilimanjaro International Airport (KIA).
The following are responses to the issues raised:
2.0 Conservation of wildlife in Tanzania
Tanzania is among the top ten countries in the world rich in biodiversity. Tanzania is also leading in wildlife conservation in Africa. It has 12 National Parks, including the famous Serengeti National Park, 34 Game Reserves and 38 Game Controlled Areas. The wildlife –protected areas cover 28% of the land surface area of Tanzania. In recognition of the good conservation works, Tanzania was awarded a conservation medal in 1995 by the Safari Club International whose headquarters is in the United States of America.
Tanzania has a number of important endangered animal species in the world. Such animal species are: Black Rhino, Wild Dog, Chimpanzee, Elephant and Crocodile (Slender Snorted Crocodile).
In 1998, the Government of Tanzania adopted a Wildlife Policy, which gives direction on conservation and advocate sustainable use of wildlife resources for the benefit of the present and future generations.
3.0 Tourist Hunting
Regulated tourist Hunting or any other type of Hunting that observes conservation ethics does not negatively affect wild animal populations. This is because Hunting ethics is based on selective Hunting and not random shooting of animals. Hunting was banned in Tanzania from 1972 to 1978. The resultant effect was increased poaching and reduced government revenue from wildlife conservation. Low revenue caused low budgetary allocations to wildlife conservation activities and the lack of working gear and equipment. When the tourist Hunting resumed Elephant population increased from 44,000 (in 1989) to 45,000 (in 1994). Elephant is a keystone species in the Hunting industry and is a good indicator in showing population status of other animal species in their habitat.
In 1989 to 1993 the government revenue from the Hunting industry increased from USD 2,422,500.00 to USD 7,377,430.00. The government earned a total of USD 9.3 Million from tourist Hunting in the year 2002. Increased revenue and keystone species such as Elephant are the results of efficient implementation of good plans and policies in conservation and sustainable use of wildlife resources.
4.0 Response to the issues raised in the article
4.1 Hunting against the law by OBC
OBC is one of the 40 Hunting companies operating in Tanzania. The Company belongs to the United Arab Emirates (UAE). Different from other Hunting companies, OBC does not conduct tourist Hunting business. The Kingdom of UAE has been the client Hunting in the Loliondo Game Controlled Area since 1992.
In conducting Hunting in Loliondo Game Controlled Area, the Company adhere to the law and regulations governing the tourist Hunting industry, namely:
4.1.1 Payment of concession fee amounting to USD 7,500.00 per Hunting block per year.
4.1.2 Requesting for a Hunting quota from the Director of Wildlife, before issuance of Hunting permit.
4.1.3 Payment of game fees as stipulated by the Government.
4.1.4 Hunting only those animals shown in the Hunting permit.
4.1.5 Contributing to the development of the Hunting block, local communities’ development projects and anti-poaching activities.
The following is what OBC has done so far:
· Contribution towards the development of the Ngorongoro District of USD 46,000.00
· Construction of Waso Primary and Secondary Schools, six bore holes and cattle dips and has purchased two buses to enhance/local transportation. Furthermore, OBC contributed TSh. 30.0M to six villages in the Hunting area, for providing secondary school education to 21 children.
· Purchased a generator and water pump worth TSh. 11.0M for provision of water to six villages. It has also constructed all weather roads and an airstrip within Loliondo area.
4.1.6. Different from the rest of the Hunting companies OBC Hunting period is very short. Normally the Hunting season lasts for six months, but OBC hunts for a maximum of four months. Few animals are shot from the Hunting permit.
4.2 Animals hunted in migratory routes.
The Government of Tanzania has permitted Hunting in the Loliondo Game Controlled Area and not in the migratory route between Masai Mara and Serengeti National Park. The Loliondo Game Controlled Area is a plain bordering the Serengeti National Park to the east.
4.3 The right for Tanzania to use wildlife in the Loliondo Game Controlled Area
The wildlife found in Tanzania is the property of the Government of Tanzania. The notion that these animals belong to Kenya is not correct. The wild animals in Loliondo Game Controlled Area do not have dual citizenship . Since some animal species move back and forth between Tanzania and Kenya it is better understood that these animals would be recognised to belong to either party during the time they are in that particular country. Animals in Masai Mara, Serengeti, Loliondo and Ngorongoro belong to one ecosystem namely, Serengeti ecosystem. However, Tanzania being a sovereign State with her own policies has the right by law to implement them. The same applies to Kenya, which has the right to implement its no-Hunting policy basing on the administration of her laws. Tanzania has therefore, not done anything wrong to undertake Hunting on her territory.
4.4 Hunting is threatening Kenyan tourism
Migratory animals move into Kenya during the rainy season. After the rainy season they move back to Tanzania. Animals that are hunted in Liliondo Game Controlled Area during this time of the year are very few. In the year 2000, only 150 animals were hunted, and in the year 2001 only 139 animals were hunted. It is therefore, not true that 80% of the animals in the border area were hunted. Based on this argument, it is also not true that Hunting conducted by OBC is threatening the Kenyan tourism industry. Tanzania does not allow Hunting of elephants 10 kilometres from the Tanzania/Kenya international boundary. (CITES meeting held at the Secretariat Offices in Geneva in 1993). This is an example of the measures taken to control what was erroneously referred to by the East African Paper as “haphazard Hunting of animals of Kenya”.
Furthermore, it is not true that the Wildlife Division does not know the number of animals that are killed. Control of Hunting is done by the Wildlife Division, District Council and other Law Enforcement agencies. The OBC does not capture and export live animals since it does not possess valid licence to do so.
4.5 OBC airplances export assorted number of live animals from Loliondo to UAE
Capture and export of live animals and birds is conducted in accordance with the Wildlife Conservation Act No. 12 of 1974 and resolutions of the Convention on International Trade of Endangered Species of Wild Fauna and Flora (CITES). The live animal trade is also conducted in accordance with the International Air Transport Association (IATA) regulations, with regard to the size of the boxes/containers allowed to transport specific animal species in order to avoid injuries or death of the same. The principle behind the live animal trade is sustainability. CITES may prohibit exportation of animals whose trade is not sustainable. On these grounds it is obvious that CITES and therefore, its 150 members recognise that the Tanzanian live animal trade is sustainable.
Live animal traders who exports animals, birds and other live specimens are obliged to adhere to the following procedure:
i) Must hold valid licence to trade on live animals.
ii) Must hold a capture permit and thereafter an ownership permit./certificate. The number of animals possessed and the number of animals listed on the ownership permit must be consistent with the number of animals that were listed in the capture permit and actually captured and certified.
iii) Must obtain an export permit for animals listed on the ownership permit/certificate.
iv) The Officer at the point of exit must certify that the animals exported are those listed on the certificate of export. The number of animals to be exported must tally with the number listed on the certificate of export.
Verification of exported animals is conducted in collaboration with the police and customs officials.
v) The plane that will carry live animals is inspected by the Dar es Salaam and Kilimanjaro Handling Companies’ Officials.
vi) For animals listed under CITES, appropriate export and import certificates are used to export the said specimens. If there is any anomaly in exporting CITES species, the importing country notifies CITES Secretariat, which in turn notifies the exporting country, and the animals in question are immediately returned to the country of export.
4.6 Other specific isues
4.6.1 Hunters are given blank permits
Companies are issued Hunting quotas before they commence Hunting activities. Each hunter is given a permit, which shows the animals that he/she is allowed to hunt depending on the quota issued and the type of safari. There are four types of safari Hunting as follows: 7, 14, 16 and 21 days safari. Each Hunting safari indicates species and numbers of animals to be hunted. When an animal is killed or wounded the officer in-charge overseeing Hunting activities signs to certify that the respective animal has been killed. If the animal has been wounded, the animal is tracked down and killed to ensure that no other animal is killed to replace the wounded animal at large. This procedure is a measure of monitoring of animals killed by hunters.
4.6.2 Good Neighbourhood Meetings between Tanzania and Kenya
There are three platforms on which Tanzania and Kenya meet to discuss conservation issues as follows:
a) The Environment and Tourism Committee of the EAC.
b) The Lusaka Agreement. In the Lusaka Agreement Meeting conservation and anti-poaching matters amongst member countries are discussed. The HQ of the Lusaka Agreement is in Nairobi.
c) Neighbourhood meeting. Experts in the contiguous conservation areas meet to discuss areas of cooperation between them, for example, in joint anti-poaching operations. Based on the regulations that govern the Hunting industry animals are not threatened by extinction since the animals that are hunted are old males for the purpose of obtaining good trophies. Trophies are attractions in this Hunting business. It is on this basis that tourist Hunting is not discussed in the said meetings, because is not an issue for both countries.
4.6.3 OBC airplanes flies directly to and from Loliondo without passing through KIA
The Tanzania Air Traffic Law requires that all airplanes land at KIA before they depart to protected areas. When the airplanes are at KIA and DIA the respective authorities conduct their duties according. The same applies when airplanes fly to UAE. They are required to land at KIA in order to go through immigration and customs checks. The allegation that OBC airplane does not land in KIA is therefore false. Furthermore, Tanzania Air Traffic Control regulates all airplanes includingly, OBC airplane at entry points.
4.6.4 OBC sprays salt in some parts of the Loliondo Game Controlled Area in order to attract animals from Serengeti National Park.
These allegations are baseless since the Tourist Hunting Regulations (2000) prohibit distribution of water and salt at the Hunting site in order to attract animals for Hunting. Besides the Game Scouts who supervise Hunting had never reported this episode. Furthermore, there are no reports that OBC is responsible for wild fires that gutters the south of the Serengeti National Park.
4.6.5 Cancellation of OBC block permit in 1999 since it was involved in the exportation of live animals.
This allegation is not true. The truth is that Hunting blocks are allocated to Hunting companies after every five years. The allocation that was done in 1995 expired in 1999. The next allocation was done in year 2000 and the companies will use the allocated blocks until 2004.
4.6.6 The UAE Royal Family contributions to the Wildlife Division
This is true. The Wildlife Division had received support from the UAE including: vehicles, transceivers and field gear in 1996. This was part of the fulfilment of the obligation by all Hunting companies to contribute towards conservation and anti-poaching activities.
Records in the Ministry of Natural Resources and Tourism show that there is no other District in Tanzania with Hunting area, other than Ngorongoro District, that receives enormous funds from Hunting business for community development. OBC contributes up to TSh. 354,967,000.00 annually for community development in Loliondo.
The Government of Tanzania has no reasons to stop the Hunting activities in Loliondo Game Controlled Area. The government sees that local communities and the Ngorongoro District Council benefit from the Hunting industry.
Edited by - lifer on 04/16/2002 20:57:41
Extract Author: Yannick Ndoinyo
Extract Date: 17 august 2002
ISSN 0856-9135; No. 00233
A rejoinder to the Ministry’s press release on Loliondo and OBC
We are replying in a critical analysis to the Ministry of Natural Resources and Tourism Press Release in the East African paper of April 1-7 2002 regarding the "Game Carnage in Tanzania Alarms Kenya" in the same paper (East Africa February 4-10 2002). Special reference was given to Hunting activities by OBC and our analysis base to same Company.
As it appears in the Press Release, OBC is the property of the UAE, but in reality some top influential people in Tanzania have some shares in the company. The OBC has been in Loliondo since 1992, even though ever since the whole local community in Loliondo refused to accept its presence and goes to the present day. OBC, to us, is not a normal Hunting company. It seems as if the Company has the right of ownership over land and other natural resources like water and wildlife. OBC has constructed expensive and luxurious houses, airstrip and big godowns on water source without the local people’s authority while they depend on such water for dry season grazing. Our surprise is that the government has always denied this fact and defended the Company. Why? The Company may be adhering to regulations and laws governing the tourist Hunting business in the books only and not practically. There are no monitoring schemes to make sure that the Company adheres to the said regulations.
It is true that OBC contributes 30 million to six villages, which is 5 million per village, and it was initially 2.5 million per village. The amount was raised two years ago. The issue here is that the amount was determined by OBC alone and therefore paid when they feel like doing, no binding mechanisms to endure payment on regular basis. The former OBC director was once quoted as saying, "I am paying this money as this money as a goodwill only because the government does not wish me to do so". The amount however does not compensate or match the resources extracted from the land of six villages. The implication is that OBC has entered into agreement with the government only and not with the villages. The provision of education to 21 children as indicated the Press release, is basically not true or correct.
The 30 million is the annual goodwill contribution from OBC to the six villages and not purposely meant for education only. The plan to utilize this money is upon the villages themselves.
It is also true that OBC has constructed Wasso secondary school but not Wasso primary school. The secondary school, which was built for the six villages in which OBC operated and the whole of Ngorongoro district has been taken by the government thus limiting the number of children hailing from these villages and Ngorongoro district an opportunity to obtain education from the same school.
In regard to bore holes, there are only four known boreholes and all these are built in Wasso and Loliondo townships. There are no any boreholes existing in any of the six villages, except only that a water pump machine, which currently does not work, was purchased for Mondorosi hamlet (Kitongoji) of Soit-Sambu village. Again, there is no virtually any cattle dip that OBC did dig or rehabilitated in six villages as mentioned in the Press Release from the Ministry of Natural Resources and Tourism. The information that OBC purchased a generator and water pump worth 11million is to six villages for water provision is false and misleading. Most villages in Loliondo have water problems and it is impossible for a generator to sustain a single village leave alone six villages. No single village has received such service from OBC.
In regard to transport, the two buses were either bought or just brought as second hand vehicles. These buses are expensive to run and spares are not easily obtained. At present they are just grounded at Ngorongoro district workshop/garage. There was a time the councilors debated whether or not to sell because of difficult management. At present the people of Loliondo use an extremely old SM bus for transport.
Again in the aspect of transport, there are no all weather roads in Loliondo that OBC constructed as it was said the Press release. Roads in most parts of Loliondo are murramed roads and mostly were constructed by Ngongoro district council using money from TANROADS and not OBC.
It is indeed true that OBC Hunting period is very short. There is a lot that can happen in short period especially if the team of hunters is composed of professional hunters. Our concern here is the interference and interruptions that OBC causes to the life systems of the people in Loliondo. The Maasai cannot resume their grazing patterns and often are forced to move by OBC. Where should we graze our cattle while our grazing land is occupied and by the Hunting company and protected by the gun? In six villages of Loliondo, five operate non-consumptive tourism that gives them more earnings except one which is dominated by OBC. The villages can now send children to school, construct basic infrastructure like health centers, classrooms, teacher houses water supply and food security ultimately eradicating poverty. This is all done using the money from the non-consumptive tourism. The OBC constantly interrupts this system and agreement between in the villages saying that the villages have no right to operate such tourism on ‘his land’. Is it his or our land? The other major problem besides the Arabs is the constant reprimand from the government, as it discourages this kind of tourism business that benefits the local people in the villages more. We favour this kind of tourism because it does not disturb our normal pattern of life system. At the same time it does not kill wild animals, they just camp and go to the Serengeti Park. The allegations that OBC airplanes fly directly from Loliondo to UAE without passing KIA and whether it exports live animals have existed and many people have spoken and written about it. However, we cannot confirm anything about without much scrutiny. We do not know much now.
Also in the Press Release the spray of salt to attract animals was referred to. The distribution of water in a certain site to attract animals was applied sometime ago. We are sure of this as it happened some years ago. What we are not exact is whether the practice continues to the present day.
In the Press Release, the records in the Ministry show that OBC pays annually 354,967,000/= for community development in Loliondo? We have some reservations in regard to such records. First of all they are just records and anything can be written. Secondly, how is it that the Ministry has such records while we in Loliondo, the base of OBC operation, do not have?
Thirdly, where is the provision in the agreement that forces OBC to annually pay to the district such amount of money? It may be that the amount is used to be paid annually but to individuals only and not to the district as it said.
In its conclusion the Ministry sees no reason to stop Hunting activities in Loliondo simply because the local community and Ngorongoro district council benefit from the Hunting business. We strongly feel that there is every reason to stop the Hunting activities in Loliondo for several reasons. First, the local community did not consent to the granting of their land to the Hunting company to the present day.
Secondly, the local community and Ngorongoro district council do not benefit in a way it should be from this Hunting business in Loliondo.
Thirdly, the presence of OBC has interrupted and interfered with our life systems including grazing, culture and other alternative means of business to the local community.
In our conclusion we feel that even though the government operates under the law set in Dar es Salaam and Dodoma without the involvement of the local people, it is very important to respect the localpeople.
Someone in the ministry who has never been to Loliondo, we firmly conclude, either wrote the Press Release, or the story was made. We suggest that the villagers or OBC people be contacted for more definite facts. Please feel free to contact us for any queries you might have regarding this article.
Tel: 0744 390 626
Extract Author: Arusha times Reporters
Extract Date: Aug 17 2002
ISSN 0856-9135; No. 00233
The Hunting plot controversy reigning within the Longido Game Controlled Area in Monduli district two weeks ago threatened the life of the American ambassador to Tanzania, Robert Royall who was Hunting in the block .
Riding in a Toyota Land Cruiser Station Wagon with registration numbers TZP 9016, owned by Bush Buck Safaris Limited, Ambassador Royall found himself being confronted by 16 armed men.
The incident took place on Saturday the 27th of July this year, at about 13.00 hours, in the Hunting block which is under the authority of Northern Hunting Enterprises Limited.
It is reported that, while Royall and his family were driving in the area, another vehicle, Land Cruiser with registration numbers TZP 3867, drove toward them and blocked them off.
Sixteen men, armed with traditional weapons including spears, machetes, doubled edged swords (simis) and clubs jumped out, ready for an attack.
However, both the ambassador, his team and driver Carlous Chalamila happened to be fully equipped and likewise drew their weapons.
Seeing modern weapons, the mob got frightened and decided to flee. But the ambassador’s driver, Chalamila followed them to find out what they wanted. Contacts were subsequently made with the wildlife department of the Ministry of Tourism and Natural Resources. Some wildlife officers were dispatched to the scene from Arusha and when they arrived, they found the attackers already gone.
Regional Police Commander for Arusha, James Kombe admitted that the incident did take place but declined comments on the issue. However, already five people have been arrested in connection with the incident, these are: Omar Mussa, David Bernard, Salimba Lekasaine, Kiruriti Ndaga and the only lady in the team, Nuria Panito Kennedy.
This week, Arusha Times learned that, the five suspects are out on bail.
Speaking by phone from Monduli, the Monduli District Commissioner (DC), Anthony Malle said there was indeed some controversy regarding the Hunting bloc of Longido Game Controlled Area (LGCA) in which two Hunting companies of Kibo Safaris and Northern Hunting Enterprises (T) Limited, were at logger heads.
Captain Malle added that, even the residents of the Singa village in the area, have been divided into two groups each supporting either companies.
The District Commissioner however, pointed out that only the Ministry will decide which of the two parties have the right to the 1,500 square kilometre Hunting bloc.
DC Malle also said that he and other district officials have already held various meetings to address the issue and together have signed an official letter which was sent to the Principal Secretary of the Ministry of Tourism and Natural Resources in order for the office to settle the matter once and for all.
Efforts to contact both Kibo Safaris and Northern Hunting Enterprises ended in vain.
The East African
Extract Author: John Mbaria
Extract Date: December 2, 2002
The East African (Nairobi) Posted to the web December 4, 2002
.. .. ..
Unlike in Kenya, the law in Tanzania promotes commercial wildlife utilisation activities such as safari Hunting and actually prohibits photographic tourism in areas declared as Hunting zones.
Under the WCA of 1974, the wildlife division can only regulate the capture, Hunting and commercial photography of wildlife.
The report adds that the director of wildlife can issue Hunting licences on village land, but he "does not have the power to give a hunter or Hunting company authority to hunt on village land without the permission of the village government."
On their part, the licensed persons are expected to seek the permission of the village government before engaging in any Hunting. However, reports indicate that the practice of safari Hunting has so far ignored this law. The report says that most Hunting companies put up facilities on village lands without the permission of the village government and the respective village assemblies.
The report gives the example of the Loliondo GCA, in Loliondo division of Ngorongoro district, where a Hunting company associated with a United Arab Emirates minister, "has built an airstrip and several large houses without the permission of the relevant village governments."
"Such actions are contrary to the VLA which, under section 17, requires any non-village organisation that intends to use any portion of the village land to apply for that land to the village council, which will then forward that application and its recommendation for approval or rejection to the Commissioner for Land."
In January, The EastAfrican published an exclusive story on the manner with which the Hunting company conducts Hunting activities in Loliondo.
.. .. ..
Extract Author: Indigenous Rights for Survival International
Page Number: b
Extract Date: 2/1/03
[click on the link to see the original MS Word document]
Indigenous Rights for Survival International
P.O. Box 13357
Dar Es Salaam.
Alternative E-mail: email@example.com
The United Republic of Tanzania
P.O. Box 9120
Dar Es Salaam.
If it pleases the Honourable President Benjamin Mkapa
Re: Stop the killing fields of Loliondo
I am a Tanzanian citizen, a strong believer in social justice. Under the same spirit I am the Co-coordinator of an informal group called Indigenous Rights for Survival International (IRSI). IRSI is a loose network of young people with an interest in public policy issues in Africa. We mainly discuss policy issues through emails communications and ultimately write articles in the press. IRSI as an entity takes no position on any of the discussed issues instead it simply stimulates, steers, and co-ordinates discussions and debates on public policy issues of members’ interest.
Mr. President, I have all along believed that you can stop the crime against humanity being inflicted upon the people of Loliondo, Ngorongoro District of Arusha Region by a no less authority than the Government of Tanzania.
Mr. President, Loliondo Division is located in Maasai ancestral lands in the northern part of Tanzania along the common border with Kenya. It borders the Ngorongoro highlands to the south, Serengeti National Park to the west, and the Maasai Mara Game Reserve in Kenya to the north. The Loliondo Game Controlled Area (LCGA) encompasses an estimated 4,000 sq km. There is no physical barrier separating the LGCA from other protected areas. It is a continuous ecosystem. LGCA was initially established in 1959 as a Game Reserve by the British colonialists under the then Fauna Conservation Ordinance, Section 302, a legal instrument the colonial authorities used to set aside portions of land for wildlife conservation. The legal status of the reserve was later changed to that of a Game Controlled Area to allow for commercial Hunting, a status that defines LGCA today and haunts its wildlife.
Mr. President, Loliondo forms an important part of the semi-annual migratory route of millions of wildebeests and other ungulates northward into the Maasai Mara Game Reserve and Amboseli National Park in Kenya between April and June, and returning southward later in the year. The survival of the Ngorongoro-Serengeti-Maasai Mara ecosystem and the wildlife it supports is linked to the existence of Loliondo and other surrounding communal Maasai lands in Tanzania and Kenya. Similarly, the survival of the Maasai people is dependent entirely upon the protection of their ancestral land for economic viability and cultural reproduction. Land to the Maasai is the foundation for their spirituality and the base for identity.
Mr. President, the people of Ngorongoro District in general and Loliondo Division in particular have suffered for a long time various established pains such as irrational grabbing of their ancestral land for “development”, tourism (consumptive and non-consumptive) and cultivation. While the people of Loliondo have lost much of their ancestral land to cultivation, the Government is evidently supporting private investors to further put Maasai pastoralists of Loliondo at a very awkward corner.
In 1992, the administration of the former president Ali Hassan Mwinyi granted the entire Loliondo Game Controlled Area (LGCA) as a Hunting concession to the Otterlo Business Corporation Ltd (OBC), a game-Hunting firm based in the United Arab Emirates (UAE). The Government issued a 10-year Hunting permit, under the controversial agreement, to the Brigadier Mohammed Abdulrahim Al-Ali, believed to be a member of the royal family of the UAE, of Abu Dhabi in the UAE who owns (OBC). The grabbed land is a birthright land of thousands of villagers of Arash, Soitsambu, Oloipiri, Ololosokwan, Loosoito and Oloirien villages of Loliondo.
Mr. President, a Parliamentary Committee was formed to probe the Loliondo Gate saga. It revoked the dirty agreement. Strangely, a similar agreement was established.
In January 2000, OBC was granted another 5-year Hunting permit in the said area. As usual, without the villagers’ consent. OBC constructed an airstrip. The villagers have been witnessing live animals being exported through the airstrip. OBC constructed structures near water sources. Hearing of the new permit, the Maasai sent a 13-men protest delegation to Dar Es Salaam in April 2000. The intention was to sort out the matter with you Mr. President. Unfortunately, they did not see you.
However, the delegation managed to hold a press conference at MAELEZO, National Information Corporation Centre. The Maasai contemplated a number of actions to be taken against both your Government and the Arab in connection with the plunder of the resources. The Maasai said that before a mass exodus of the Maasai to Kenya the first thing was to eliminate wild animals. Thereafter, the delegation retreated to Loliondo, as gravely frustrated as before.
The general election was scheduled for 2000, so the saga had to be explained away. The official statement was that power hungry opposition politicians were pushing the elders and that all the claims by the Maasai were “unfounded” and “baseless.” To its credit, The Guardian went to Loliondo. It reported the following:
Maasai elders in Loliondo, Arusha Region, who recently declared a land dispute against OBC Ltd, a foreign game-Hunting firm, have accused some top Government officials of corrupt practices, saying the conflict is not political. The Arusha Regional Commissioner, Daniel ole Njoolay, recently described the simmering land dispute between the Maasai pastoralists and OBC, as a political issue.
Francis Shomet [the former Chairman for Ngorongoro District Council] claimed that Njoolay had misled Tanzanians to believe that the allegations recently raised by Maasai elders were unfounded and baseless. Fidelis Kashe, Ngorongoro District Council Chairman maintained, “We cannot stand idle to see our land being taken away by Arabs. We will kill all the animals in the area as these are the ones attracting the Arabs into our land” (The Guardian May 30, 2000).
The next morning Government officials were reported to have said the following:
The Minister for Natural Resources and Tourism, Zakia Meghji, yesterday assured Ngorongoro residents that no land has been sold or grabbed by Arabs in Loliondo. Flanked by the Arusha Regional Commissioner, Daniel ole Njoolay and the Director of Wildlife, Emanuel Severre, Meghji commented, “There is no clause on the sale of land in the contract signed between OBC and the six villages of Ololosokwan, Arash, Maaloni, Oloirien, Oloipiri and Soitsambu.”
However an inquiry conducted by The Guardian in Loliondo last week established that the Maasai elders were not involved in the re-lease of the Hunting block to the company. According to Megji, her probe established that the building has been constructed about 400 metres from the water source, 200 metres more than the distance recommended by law. But The Guardian investigation shows that the structures are less than 50 metres from a spring. And another spring has dried up (The Guardian May 31, 2000).
Mr. President, underline two points. First, the Minister said the building has been constructed 400 metres from the water source. Second, “The Guardian investigation shows that the structures are less than 50 metres from a spring.” Now unless one’s mathematics teacher at school was daft, there is a huge different between 50 and 400! When did 50 metric metres turn to mean 400 metric metres? Can it be claimed that the Maasai were party to this so-called agreement? I am at a loss why this-well known-Minister has not been made to face the full force of the law.
In the proposal, Brigadier Al Ali outlined the benefits of his operations in Loliondo to the Government, local communities, and wildlife conservation in the Serengeti-Maasai Mara-Ngorongoro ecosystem. Among its important objectives were:
• To conserve an area contiguous to the Serengeti National Park, which is essential to the long-term survival of the ecosystem and its migration.
• To develop a new role and image for the Arab world as regards wildlife conservation, management, and human development.
• To improve locals’ revenue, development facilities, and create employment.
• To generate revenues for the Central and District Governments.
The OBC now stands accused of self-contradiction and violation of legal and moral obligations in virtually all the above areas, resulting instead in environmental destruction; unfulfilled promises and exploitation of the local communities; and direct undermining of the stability of the region’s wildlife and natural habitats.
It has become evident that OBC had a long-term agenda for exploiting the high concentration of wildlife in Loliondo. Its Hunting operations are guaranteed by the continuous flow of wildlife from the Serengeti, Ngorongoro, Maasai Mara, and other areas. According to the International Union for Conservation of Nature, OBC "was taking advantage of migratory patterns of wildlife coming out of Serengeti."
Mr. President, be informed that the villages in and adjacent to protected areas in Tanzania have no Government-supported infrastructures. Take Ngorongoro District for instance. There is no Government hospital in Ngorongoro. It may take a week to travel from Arusha to Loliondo, just less than 400 km, depending on weather, for there is no road. There is no even a single Government advanced level secondary education school in six (repeat six) Districts in the Greater Serengeti Region. This situation brings to question the legitimacy of wildlife conservation vis-à-vis the right of rural people to lead a decent life given nature endowment in their localities.
Mr. President, the Maasai of Loliondo have for a long time accused OBC of grave human rights abuses. They have described acts of intimidation, harassment, arbitrary arrest and detention, and even torture by OBC staff, Tanzanian police and military in the name of OBC; brazen violations of grazing and land rights; and wanton environmental destruction and imminent extermination of wildlife. They have seen leaders who once opposed OBC’s practices corrupted and bought-off.
The OBC operates like a separate arm of the Government. Many people in Loliondo believe that OBC is even more powerful than the Government. The Maa word for "the Arab", Olarrabui, is often used to refer Brigadier Al Ali, and by extension OBC. The word Olarrabui has become synonymous with power, authority, brutality, fear, and entities larger than life.
Mr. President, you do not need to be a rocket scientist to comprehend that this is the clearest case of abuse of office. It is suggested, for those willing to avert disaster, the Tanzania Government included, that immediate steps be taken to put to an end the violation of fundamental human rights in Ngorongoro. As to lands lost in Loliondo, the Government is advised to return this to its owners. Land should not be grabbed senselessly. The Government, should at once, re-look into the whole matter.
Navaya ole Ndaskoi.
- The International Court of Justice
- The United Nations High Commission for Human Rights
- The United Nations Working Group on Indigenous Populations
- Human Rights Groups around the World
- Faculty of Law of the University of Dar Es Salaam
- Local and International Conservation Agencies
- Ministry of Tourism and Natural Resources
- The Attorney General
- The Chief Justice
- The Speaker of the United Republic of Tanzania Parliament
- The Press, print and electronic
- Political parties in Tanzania
- Tanganyika Law Society
- Other interested parties.
Navaya ole Ndaskoi
see also Extract 3734
The Maasai protest delegation holding a press conference in Dar Es Salaam in 2000 |
Establishing awareness of intellectual property amongst staff of your company is essential for early maximizing the value of your intellectual property and the wealth of your business and reducing the possibilities of accidental non-confidential disclosures, that could prejudice successful patent applications and negatively affect the value of your intellectual property and ultimately the wealth of your business.
Regular training sessions of staff on intellectual property are key and should include the following:
how to identify and protect intellectual property;
how to use patents to improvements of technology;
understanding Patent Process;
how to deal with confidential information (see some examples in the scenarios below);
record keeping of intellectual property, including laboratory notebooks and policy on intellectual property; and
who to contact in case of need.
The record keeping procedures and manuals will address the following questions:
has the inventor kept the idea confidential?
is there a written description of the idea and has been kept safe and confidential?
how the idea has been generated? If during a collaborative programme, then was it agreed beforehand who owns what?
is the idea a new product, a new material, a new process for making something? If so, is it patentable or protectable in any other way?
is the idea a variation in a product or material or process? If so, it is still likely to be patentable or protectable in any other way?
who generated the idea? The answer to this question is very important in the event self-employed or other third party consultants are involved in any research and development or collaborative project.
The main object of this record keeping is to track, protect and maintain all relevant intellectual property rights of the business so that intellectual property can be licensed, assigned or exploited to the fullest extent and benefit of the company.
The record keeping procedures should also include a form upon which, potential inventions should be recorded identifying the following:
Who: department and research area;
Named individuals: inventors and authors;
What: technical description;
Why: perceived novelty;
How to use the information: potential applications/markets;
What else is needed: background or third party intellectual property and information.
The correct use of laboratory notebooks by staff is also essential. In the event of a dispute laboratory notebooks may be required to be presented as legal evidence.
It is therefore recommended that:
permanent bindings are used on notebooks (loose leaf books should be avoided to prevent possible removal or substitution of pages);
pages should be numbered and any additional drawings cards or computer printouts should be permanently attached to the notebook clearly identified and have reference made to them in the notebook;
all projects related and other activities, such as breaks in research due to secondment or holiday should be recorded factually; and
the notebook should be reviewed regularly by someone who understands the technology involved, each page should ideally be signed by a witness and again the choice of the witness is important and should not be someone who may be nominated as core inventor. The witness should also sign and date and graphs, chart, printouts, which are inserted into the laboratory notebook.
In addition, to the use of appropriate record keeping procedures and notebooks, an evaluation of IPR policy should be adopted. Such an evaluation should include factors, such as potential market, market, impact, competitive products, timing, intellectual property protection available and experience in the field concerned.
Finally, once relevant intellectual property rights have been identified, protected, exploited and enforced, it is advisable that, regular audit of such rights is undertaken to ensure that, the intellectual property rights reflect the current needs of the business and that expenditure is limited accordingly.
How to take care of confidential information:
Scenario 1 –
Have you been asked to sign a confidentiality undertaking? If so, please check that it is only confidentiality restriction and not a transfer of intellectual rights.
Obtain express confirmation from the discloser, that the information is not confidential, where possible, before disclosure.
Make a written record of what was disclosed, by whom and when.
Please remember that, an obligation, to keep information confidential, includes the obligation of not disclosure or not use of the information, without the permission of the person to whom the obligation is owned.
Scenario 2 –
Put in writing or some other permanent form.
Mark any documents with appropriate confidentiality and IPR disclaimers.
Keep a copy of what is disclosed and a record of when and to whom.
If an oral disclosure is made in confidence, confirm in writing what was disclosed and what was given in confidence.
Have the recipient sign a confidentiality undertaking of the disclosure.
Scenario 3 –
Consider whether anything in the paper describes a new device, chemical compound or manufacturing process or a significant improvement or modification to any such matters.
Do not disclose anything, without first considering the possibility of the content of the papers being patentable in whole or in part.
Consider whether there are any restrains, under any relevant agreements (including research and development or collaboration agreements).
Keep an eye on any relevant timetable for confirming publication.
Request that, the publisher confirms confidentiality on receipt of paper pending decision on publication.
Please, always remember that, any document exchanged, should be clearly marked as being confidential.
Scenario 4 –
Consider what background IPR, if any, are free from obligations of confidentiality and may be introduced to the project.
Prior to disclosing any information to third parties, have a confidentiality agreement signed. Such agreements, may take many forms and the terms should be adjusted in accordance with the particular circumstances.
You should always include the following:
- identification of parties;
- what information is to be kept secret; and
- for how long.
If in any doubt, consult your legal adviser.
The World Intellectual Property Rights Organization, WIPO has organized a three day workshop to educate the public on intellectual property rights. The initiative started today at the Coco palm hotel. The drive which is aimed at sensitizing St. Lucians to intellectual property rights issues will also include a public forum for persons interested in learning more about the issue.
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The expression, ‘intellectual property’ has come to be internationally recognized as covering mainly two branches, namely; ‘industrial property’ and ‘copyright’. Patents, industrial designs and trade marks used to be considered as different kinds of industrial property. In Bangladesh (during the time of the then British regime), the first legislation of its kind, on copyright was introduced in 1914, which was mainly based on the British Copyright law of 1911. After the independence from Britain new law on copyright was promulgated in 1962. The Copyright Ordinance 1962 has been replaced by the new copyright act of 2000. Now in our country, Copyright law is regulated by the Copyright Act 2000. This is done because of the prevailing situation in Bangladesh and around the world.
In case of patent and designs we have law which we inherited from our Colonial ruler. Patent rights are created by statute and governed by the Patent and Designs Act 1911. We have also law related to trademarks and it is regulated by the Trademarks Act 1940.
Intellectual property has acquired an internationally recognized character. Now it is regarded as “one of the most important sectors” of international law, having its source in various international conventions. At present, each and every country is trying to shape or reshape their legislature, relating to intellectual property; in the light of those international conventions.
The convention establishing the World Intellectual Property Organization (WIPO), concluded in Stockholm on July 14th, 1967, provides that, ‘intellectual property’ shall include rights relating to:
1. literary, artistic and scientific works,
2. performances of performing artists, phonograms and broadcasts,
3. inventions in all fields of human endeavor,
4. scientific discoveries,
5. industrial designs,
6. trade marks, service marks and commercial names and designations,
7. protection against unfair competition,
And all other rights resulting from intellectual activity in the industrial, scientific, literary or artistic fields. (Article 2/VIII) 1
Intellectual property protects application of ideas and information that are of commercial value.2
By virtue of a number of international conventions such as the, Berne Convention and the Universal Copyright Convention, copyright acquired in one country extends to other countries which are member of these conventions. Other intellectual properties are beginning to acquire the same nature as well.
The principles of intellectual property law are substantially the same in all countries with little variation to meet the national requirements of each of the countries.
Bangladesh, as a least developed country, has also enacted intellectual property law in its national legislature.
The Trade Related aspects of Intellectual Property rights (TRIPs) were included as integral part of World Trade Organization (WTO) due to pressure and interest of basically transnational companies and developed countries to ensure maximum profit or interest out of intellectual property in international trade.3 The TRIPs Agreement has been called the most ambitious international intellectual property convention ever attempted. TRIPs agreement has established the protection of intellectual property as a major part of the multinational trading system embodied in the WTO. As one commentator observer, intellectual property is now a key component of this trading system, “the protection of intellectual property is one of the three pillars of the WTO”.4
Bangladesh is now enjoying the transitional period that has been fixed by WTO under the light of TRIPs. So, we do not have a clear idea about the relation between intellectual property rights and economical development. But, it is certain that after the expiration of transitional period, Bangladesh will have to face a serious and severe situation because of the weak legal framework relating to intellectual property.
Now, in brief, the main features of intellectual property laws are given below:
1. The Patents and Designs Acts, 1911:-
(Act no. 11 of 1911)
The salient features of our existing Act are as follows;
Part-1; of the Patents and Designs Act, 1911; is about of patent. This part starts from section-1 and extends to section-42.
Section 2(8) contains the definition of invention; it proceeds as – invention means any manner of new manufacture and includes an improvement and all alleged invention. In sub-section (10) of the same section, provides the explanation of manufacture, “manufacture includes any art, process or manner or producing, preparing or making any article and also any article prepared or produced by manufacturer.”
Section-3 of the Act provides about the manner and mode of an application for patent. According to section-3(1); a patent application can be made by any person, whether he is a citizen of Bangladesh or not. An application can be made, alone or jointly with any other person.
According to section-4 of the Act; there are some provisions about specification. Wection-4(2) states that, a complete specification must particularly describe and ascertain the nature of the invention and the manner in which the invention is to be performed.
Section-10 of the Act, discusses about the topic, “grant and sealing of patent”. According to section-10(1); if there is no opposition of the patent application or, in case of opposition, if the determination is in favor of the grant of a patent, a patent shall, on payment of the prescribed fee, be granted, subject to such conditions if any as the Government thinks expedient, to the applicant, or in the case of a joint application to the applicants jointly, and the Controller shall cause the patent to be sealed with the seal of the Patent Office.
Term of patent, is laid down in section-14(1) of the Act. The duration of patent is sixteen years from its date. Section 15 of the Act; also provides rules regarding extension of the term of patent.
If any patent has been ceased, owing to the failure of patentee to pay any prescribed fee within the prescribed time, the patentee may apply to the controller in the prescribed manner for an order for the restoration of the patent. Section-16 of the Act deals with the matter, restoration of lapsed patent.
Section-22 of the Act contains provisions of “Compulsory Licenses and Revocation”. Both the Govt. and High Court Division are empowered to grant compulsory license or revocation of patent. Any person interested may present a petition to the Govt., alleging that the demand for a patented article in Bangladesh is not being met to an adequate extent and on reasonable terms and praying for the grant of a compulsory license, or in the alternative, for the revocation of the patent.
According to section-25 of the Act, if any patented invention or the mode which it is exercised; is mischievous to the state or prejudicial to the public; Govt. can declare the patent, revoked, by notification in the official Gazette.
Section-26 of the Act; provides the grounds, on which; a patent can be revoked by the High Court Division. Section-26(2) also provides, who can present a petition for revocation of a patent.
Section-30 of the Act; declares that; an innocent infringer of a patent; is exempted from liability or damages for infringement.
Part-II of the Act; describes all about ‘Designs’ or industrial designs. According to section-43(1); any person claiming to be the proprietor of any new or original design not previously published in Bangladesh, can apply to the Controller for the registration of that particular design.
According to section-47(1); the proprietor of a registered design, shall subject to the provisions of this Act, have copyright in the design during five years from the date of registration.
Section-51(A) of the Act; narrates that, any person interested may present a petition for the cancellation of the registration of a design. Such petition should be presented to the High Court Division.
Part-III of the Act; is about, Patent Office and proceedings there at.
According to section-59 of the Act; every register kept under this Act shall at any convenient times be open to the inspection of the public, subject to provisions of this Act.
Section-65 of the Act; states that, in any proceeding under this Act, the Controller shall have the powers of a Civil Court for the purpose of receiving evidence, administrating oaths, enforcing the attendance of witness compelling the discovery and production of documents, issuing commissions, for the examining of witnesses and awarding costs and such award shall be executable in any Court having jurisdiction as if it were a decree of that Court.5
2. The Trade Marks Act; 1940:
(Act No. V of 1940)
The basic principles of the Trade marks Act are described below:
Section-2 of the Trade marks Act contains definitions that are related to this Act. As to, section-2(k); trade mark means a mark used or proposed to be used in relation to goods for the purpose of indicating or so as to indicate a connection in the course of trade between the goods and some person having the right, either as proprietor or as registered user, to use the mark whether with or without any indication of that identity of that person.
The establishment of Trade marks Registry at Dhaka, appointment of the Registrar and Deputy Registrar are laid down in section-4 of the Act.
According to section-5 of the Act; the registration of a trade mark, requires distinctiveness. Purpose of such distinctiveness is to distinguish those particular goods from the others, which have similarity in nature.
Any mark, containing scandalous design; or be likely to hurt the religious susceptibilities of any class of the citizens, or to be contrary to any law for the time being in force or morality is prohibited for registration of that mark. Section-8 of the Act, say so.
Section-16 of the Act; provides that; when an application for registration of a trade mark has been accepted and either has not been opposed or having been opposed, has been decided in favor of the applicant, the Registrar shall, registers the said trade mark.
Section-18 of the Act, says that; the registration of a trade mark shall be a period of seven years, and may be renewed from time to time in accordance with the provisions of this section.
According to section-20 of the Act; no person shall be entitled to institute any proceedings to prevent, or recover damages for the infringement of an unregistered trade mark.
According to section-46 of the Act; any person aggrieved can apply to the High Court Division or the Registrar, for the cancellation or verification of the registration of a trade mark on the ground of any contravention of, or failure to observe a condition entered on the register in relation there to.
Chapter-IX of the Act; is specially written down for textile goods.
Chapter-X of the Act, states the provisions regarding offences and restriction of use of Royal Arms and state emblems.
According to section-73 of the Act, any suit for the infringement of a trade mark or otherwise relating to any right in the trade mark; shall be instituted to a District Court having jurisdiction to try the suit.6
3. THE COPYRIGHT ACT; 2000:
(Act No. 28 of 2000)
A short overview:-
Section-2 of the Copyright Act; provides all the definitions related to copyright and so on.
According to Chapter II; section -9, 10 and 11; Copyright Board will consist and the post of Register has made. The board is a quasi-judicial body; while working, it would be deemed as a Civil Court.
Definition of copyright is laid down in Chapter III; section 14 of the Act. Section 14(2) includes, ‘computer programs’ as a subject to this Act.
Chapter IV; deals with the ownership of copyright and the rights of the copyright owner. This Chapter starts from section-17; ends to section-23.
Chapter V of the Act; describes all about the term of different types of copyright. Generally the term extends from the lifetime of the author until sixty years from the beginning of the calendar year next following the year in which the author dies; is at section-22. The section also provides that, copyright shall subsist in any literary, dramatic, musical or artistic work (other than a photograph). Section-25 to section-32 of the Act narrates the duration of copyright for different types. For instance, posthumous work, cinematograph films, sound recordings, photograph, anonymous and pseudonymous works, Govt. works, work of any local body, work of international organizations.
Section-50 of the Act; deals with compulsory licenses in works withheld from public.
Chapter X of the Act is about, registration of copyright.
Chapter XII; section-71, describes about infringement of copyright and section-72 of the Act; provides acts not to be infringement of copyright.
Chapter XIV; is about Civil remedies, that are available against infringement of copyright. Exclusively, section-76 of the Act; provides, remedies against such infringement. According to section-81 if the Act; the Court of District Judge, is the Court of first instance of such proceedings.
Chapter XV of this Act; describes offences and punishment. According to section-82 of the Act; any person who knowingly infringes or abets the infringement of copyright; shall be punishable with imprisonment for a term which shall not be less than six months. But, which may extend to four years and fine which shall not be less than fifty thousand taka but which may extend to Taka two lacks. But; where the infringement has not been made for gain in the course of trade or business; the Court may, adequate and special reasons to be mentioned in the judgment, impose a sentence of imprisonment for a term of less than six months or a fine of less than fifty thousand taka.
According to section-95 of the Act; an order made by the Registrar is a subject to appeal to Copyright Board. According to section-96; against any order; made by the Board; any aggrieved person, within thirty days of such order can file a petition of appeal to the High Court Division.7
PROBLEMES AND INSUFFICIENCY OF THE EXISTING INTELLECTUAL PROPERTY LAWS:
PATENTS AND DESIGNS ACT:-
The defects are;-
Like any other national patent system, novelty is an essential requirement of an invention to be patentable under the existing law of our country. But it is clear from the definition of ‘invention’ in section-2(8) of the 1911 Act; the invention to be patentable need not have a commercial pecuniary success. In other words, utility of an invention is not required by this definition. Although lack of utility is a ground for revocation of a patent under section-26(f) of the Act, till now this has not been included in the definition of invention. The 1911 Act makes no difference between patentable and non patentable inventions. Considering the public interest certain items should be kept outside the domain of patent protection. Since under the existing law no item has been excluded from the domain of patent protection, any type of new invention may obtain patent protection although it should not be patentable for the greater interest of the public at large. Under the existing law, the application for a patent must contain a declaration that the applicant is the true and first inventor or the assignor or legal representative of such inventor. But the term, “true and first inventor” is left undefined. Thus any person importing an invention for the first time into Bangladesh or any person to whom an invention is first communicated from outside Bangladesh can claim to be the “true and first inventor” of that invention and thus, can take the advantage of the vagueness of the term. Here; a fact should be added that, such practice already has begun. Some pharmaceutical companies don’t having any research and develop cell, importing some drugs, and by claiming patent, under the ambiguity of the Act, getting patent. A statistic says that, at the 1st and 2nd week of March, 2006; on an average, 40 patents were granted by the patent office, in each week. Thus, the local producers are forced to remain non productive on that particular patented items. The standard of examination varies from country to country. In some countries like Netherlands, Germany, U.S.A. and Japan it is rigorously involving an extensive search for both novelty and obviousness among documents published in many countries, over a period of many years. But according to our existing patent law; examination is less rigorous involving for novelty only, and the extent of search is restricted. Term of patent protection shall, as laid down in section-14, be sixteen years from its date. But the term is not sufficient enough for the exploitation of the patent right. What constitutes infringement of patent isn’t defined in the Patents and Designs Act, 1911. Section-29(1) of the Act only says that, the patentee has a right to sue against the infringer during the continuance of the patent acquired by him. Section-30 of the Act says that, a patentee shall not be entitled to damages against an innocent infringer. But, ignorance of law is no defense. The infringer should pay for the project as compensation of damages. Under the existing law any process or manner of producing, preparing or making an article in patentable as appears from section-2 (8) read with section-2(16). But, the 1911 Act does not confer upon the patentee the exclusive right to exercise the process. The 1911 Act contains provisions regarding compulsory licenses of patent rights but the terms and conditions of conditions of compulsory licenses are not detailed in the Act.
TRADE MARKS ACT:-
1. Section-22(3) of Trade Marks Act defined infringement in a very narrow sense though the act has provisions against infringement.
2. Although Chapter X of the Act, describes about offences and restrictions of use of Royal Arms and state emblems but this chapter does not extends to the infringement of any registered trademark. Any deceptive use of any registered trademark; is not also included in this chapter.
3. There is no provision for protection of internationally recognized trademark in our existing Trademark Act 1940.
Our existing copyright law has been enacted in line with the copyright law of India. It has been enacted to cope with the prevailing international set up of copyright system. So, preventive measures have been adopted to tackle the future complications in copyright sector.
NECESSARY PROPOSAL FOR REFORMATIONS: The term ‘intellectual property’ is still at its nascent stage in our country and people are not aware of the concept and importance of intellectual property. But, in international arena, the concept and coverage of intellectual property is growing so fast than any other brunches of law. We have intellectual property laws but these laws are not sufficient to tackle the challenges that are imminent and threatening us. Keeping in mind the Trade Related Aspects of Intellectual Property Rights (TRIPs) agreement and other relevant conventions the following reformation proposals can be made:
Reformations of the Patents and Designs act:
ü The Patent and Designs Act, 1911 should be revised thoroughly.
ü Essential requirements of patentable invention should be described clearly and there must be a clear distinction between patentable and non patentable inventions.
ü The standard of examining a patent application should be made more effective.
ü Term of patent protection shall, as laid down in section 14, be 16 years from its date. But the term should be extended to 20 years for patent and the term of a design, according to sec – 47(1) is 5 years from the date of registration, should be extended to 10 years.
ü Existing Act does not have any definition of infringement, it should be included.
ü The part of “designs;” have some confusing words, as in that part the term, “copyright” has frequently used. But it may create confusion. Such confusion should be effaced.
ü The Act does not have anything to do about the protection “Geographical indication” which could result a huge loss of losing our culture & heritage. So, it should be included.
ü Provisions relevant to PARIS convention should be incorporated.
ü The administrative provisions and complications should be avoided. The complications should be made more subject to judicial decision.
ü The provisions of offences and penalties should be revised and reformed with the need of the time.
Proposal for reformations of The Trademarks Act:
After studying present Trademark Act and different conventions, it is clear that our existing Trademarks Act 1940 should be amended as well. The following suggestions can be made:
ü The reformed trademark law should have a wide view about the marks which can be registerable and which marks cannot be.
ü Infringement of trademark should be defined more accurately. Besides, punishment for infringement should be made stricter.
ü How will we protect our renowned trademarks in international market and reciprocally how other countries trade marks can enjoy protection in our local market, should be defined in our trademarks act, with an assertive view.
ü Offences and penalties, in respect of violation of any provisions of this act, should be made more effective and harder.
ü If any complication or confusion arises in practicing of the act, the judicial body should be involved with more vigilance.
Proposed reformations to the copyright law:
The Copyright Act 2000 has been enacted to cover the rules and to cope with the international copyright system. Our existing copyright law has been enacted in line with the copyright law of India. It has been enacted to cope with the prevailing international set up of copyright system. So, preventive measures have been adopted to tackle the future complications in copyright sector. From the face of the Copyright Act 2000, it seems that our copyright law has fulfilled the need of the time.
Though computer programs, tables complications including data base are recognized to have the copyright protection, there is no legal recognition for transaction carried out by means of electronic data or other means ‘e-commerce’ which involves the use of alternativeness to paper based method of communication and storage of information to facilitate filing of documents with government agencies .
The growing global importance of the cyber law is posing new challenges and in view of the peculiarity involved in the fields, the understanding between the nations of the world by treaties or covenants, may be of considerable importance in the absence of which the implementation of the legislation would be near to impossibility. Our present act should be amended to fulfill the shortage.
Finally we can hope that a stronger protection system of the intellectual property rights, a qualified commission to observe the protection and thus policy making options to bring civil remedies for the violation of the rights and finally a complete law regulation in all sectors of intellectual property rights will surely lead us to a better future.
Notes and References:
Background reading materials on intellectual property: WIPO (Pp-3, 4).
P.Narayanan.(pp-1). TRIPs agreement. John Madely hungry for trudeyzed books UK (2000) pp-96-97. The Patents and Designs Act 1911 The Trade Marks Act 1940 The Copyright Act 2000
Most people are attentive of the numerous benefits of owning a trademark registration. Trademark registration in India becoming familiar with complete customer satisfaction. Trademark registration is the protection agreed by the government to the business entities as to reduce the possibility of getting the advantage of the business by others by the way of misuse and to raise the opportunities keeping the mark exclusive under the eye of law.
Generally, brand registration refers to the trade mark used to discriminate the goods or services among the consumers. The business group sells their services or goods under the precise name or brand that is called trade mark. Therefore, the brand is registered in order to evade the repetition or use the same mark by others. In vision of this, the brand registration referred to as trademark registration. Trademark brand was initially developed as a name, term, design, and symbol. Powerful brand can bring success in bloodthirsty and financial markets and thus become the markets worthless assets.
Trademark brand equity dealings the value of brand to the trademark owner. The brand name is used interchangeably with brand to designate written or spoken linguistic rudiments of the brand. Brand name is a form of trademark which identifies the brand owner as the money-making source of products or services. The brand owner may ask for to protect the proprietary rights in relation to a brand name during . Trademark brand is a appliance to create monopoly so that the brand owner can obtain some of the reimbursement to those related to decline price competition. There is legal magnitude as it is essential that the brand names and trademarks are protected by all means. An existing brand name can be used as a vehicle for new and modified products. Individual brand names allow greater suppleness by permitting different products to be sold without puzzling the consumer.
The trademark is registered for the business name, brand name and logo as to discriminate, popularize, create the goodwill and put aside the mark from competitors and fraudulent. The trademark office is an organization to provide protection to the inventors and dealing for their inventions and trademark registration in India provide protection and intellectual property recognition.
In addition, if some business entity desires to extend its dealing in more than one or several countries, it can ensue with International Trade mark registration. The titled name International Brand Registration is the usual form of the International Trade mark registration. It is meant, when the registration is done through any International pact, that gives the protection in all the countries allied with the treaty.
The is an agency, which provides protection to the inventors and business for their inventions and trademark registration for the product and intellectual property identification. The office is provided with funds by the fees, which are charged for processing the patents and trademark. The applications to trademark registration are examined by the trademark office.
Master sf Business Administratration Banking and Finance (MBF): Integrative Courses: MS-94 Technology Mangement
Video Rating: 0 / 5
- the various types of merchandising rights including intellectual property rights and monopoly in commonplace articles
- use of well-know images and features in advertising and unfair competition
- use of registered trade marks and service marks
- use of copyright and design law
- the civil remedies available
- international framework for the protection of merchandising rights
Find More Intellectual Property Rights Products
The service industry is a tricky landscape, due to the fact that, often times, no tangible property is transferred back and forth. Sure, you may purchase a software program that comes on a disk, but in reality, you are merely purchasing the right to use that program, not the rights to the program itself. The service industry is a world dominated by intellectual property, and protecting yours should be your utmost goal.
First you must define if intellectual property is yours or not. Let’s say you hire an employee to create a software program for you. He creates it, puts his name on it, and hands it over to you. If left unchecked, both of you could walk away thinking that you own the rights to the software, but in reality, only the employer does. Other items may fall under the umbrella of “intellectual property” including trade secrets, client information, and more. Your employees might have access to that information, but in order to legally prevent them from using that information to their own benefit, you must first get them to sign an employment contract informing them that they are not allowed to distribute any information gathered while under your employ.
As an employer, you should take pride in the knowledge, contacts, and experience you have accumulated over your years of doing business. Don’t risk letting that information out by avoiding employment contracts or other intellectual property agreements. Furthermore, don’t let subcontractors walk away with information you paid for, just because they created it. Understanding your rights as an employer is the first step towards protecting them. The final step is having the foresight to create contracts which will legally protect you should the situation ever call for it.
As part of the Center for Freedom and Prosperity Foundation’s video project, we sponsored a contest for students at the University of South Florida. The winning entry discusses the role of intellectual property rights.
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The relevant legal norms from the above point of view, although they involved Intellectual property rights The exercise of acts done some direct or indirect restrictions, but these provisions are very fragmented, incomplete, unclear, not specifically from the perspective of preventing the abuse of intellectual property rights to make rules; the majority of existing legal norms applicable to the foreign economic Trading Activities, and not generally applicable to the Chinese market, the exercise of acts of intellectual property, so its very limited scope; from the content point of view, it also needs to be common practice in the world at present updated and perfect.
Present, China has no clear regulation of the exercise of acts of intellectual property ” Antitrust Law “, Also led to the present does not exist in China, the corresponding administrative law enforcement and judicial practice, similar Microsoft , Cisco And DVD Alliance patentee of intellectual property rights abuse in China still can not be effective monopoly regulation.
Order to achieve Competition Request on behalf of the wider interests of society more important, China has continued to strengthen intellectual property protection, but also abuse of intellectual property should be to rule them. But now, people are more concerned about the previous problems, while the latter issue also seriously enough. Based on this, the urgent need to promptly establish and improve regulation of intellectual property rights abuse in China antitrust law system, give full play to intellectual property rights in the legal system to encourage innovation and promote scientific and technological progress of the positive role of the city and to prevent the abuse of intellectual property rights to maintain a free and fair economic competition order.
The “anti-monopoly law,” the exercise of intellectual property regulation act to protect the interests of Chinese enterprises is significant. Trend of economic globalization is bound to be more Chinese enterprises to the market, they may encounter in the domestic market monopoly of multinational enterprises do nothing, and their behavior in foreign markets are subject to strict Competition Law problems. For example, earlier this year, some U.S. Pharmacy Enterprises in China 4 pharmaceutical companies in the United States filed an antitrust lawsuit accusing the Chinese companies in the U.S. market on price cartel. Can be expected, as more and more Chinese enterprises enter the international market, Chinese enterprises in foreign countries have been accused of anti-monopoly situation will be more. Therefore, whether for the maintenance of domestic order, or free and fair competition in international economic exchanges in safeguarding their own interests, China should be established as soon as possible “anti-monopoly law.” The core is to balance and handle competition in intellectual property protection and maintenance of the conflict between the ultimate goal of full respect for and protection of intellectual property rights, encourage innovation and play Excitation The role of competition, but also effectively prevent illegal monopoly by improper use (ie abuse), so on behalf of the interests of society as a whole will not be free and fair competition order of destruction; both the full protection of market competition, but also realistic and reasonable for the time being take care to limit competition business needs, a reasonable balance of intellectual property transactions, the parties (developers, producers, consumers, etc.) interests, in order to promote China’s scientific and cultural innovation and the parallel development of economic competition.
Specialized intellectual property law in the regulation of intellectual property rights abuse in the sound system, to further clarify, refine the terms of the abuse of intellectual property, which makes the related intellectual property rights infringement proceedings for the alleged infringer to provide a clear defense based on, or to the right people can be counter-claim, or even be prosecuted separately. This also requires our country, “Patent Law,” “Trademark Law”, “Copyright Law” and other specialized intellectual property law further modified, or other supporting measures. For example, should further improve our procedures for intellectual property litigation related to the legal system, clear that abuse complaint against compensation; in intellectual property cases should be taken provisional measures before litigation, scrutiny, careful decisions; improve the institute confirmed that the provisions of such non-infringement lawsuit .
In addition, to effectively safeguard the authority of our laws, to be taken seriously TNC Intellectual property dispute resolution and political trend.
As intellectual property rights are private rights, in essence, is a civil right, then the intellectual property dispute resolution, like other civil rights, are primarily to the court through the trial process carried out. Related Intellectual Property Rights in China Laws and regulations After recent amendments, has been full compliance with WTO rules, the minimum requirements, in some ways even more than that can provide effective intellectual property legal relief channel.
However, some multinational companies and Chinese enterprises in intellectual property disputes occur, often by virtue of its strong bargaining power directly to the Chinese government or by their home governments to pressure the Chinese government to meet their special benefit The purpose of treatment, so that would could also go through legal channels, in accordance with strict legal procedures to resolve political. The Chinese Government has always considered positive for a variety of administrative resources to address the use of foreign-related intellectual property disputes. Tasted the sweetness of the multinational corporations, the more value this way, and attract more multinational companies to make the same choice. Moreover, under WTO rules, the MFN principle, a country our government to give preferential treatment for enterprises in this regard should also be automatically and unconditionally to other WTO Organization Members of the business. This is both an expression of some multinational companies in China to protect intellectual property rights and the rule of law framework for the establishment of disrespect and distrust, often the interests of the enterprises in China causing undue damage.
Is therefore recommended that our government should be treated with a number of multinational corporations will solve the political disputes of intellectual property rights of the tendency to be guided in accordance with our existing law, to follow strict legal procedures to solve problems, in particular, be careful not to as being a political Pressure The expense of the legitimate interests of Chinese enterprises. Access to the judicial process for foreign-related intellectual property disputes, our courts must strictly in accordance with Chinese laws and relevant international rules, rule out all kinds of interference imposed by foreign parties, according to a just ruling, effectively protect the legitimate rights and interests of Chinese enterprises. Which claims the rights for the source of the right to conduct a strict review of the situation, to prevent the use of intellectual property claims on behalf of persons engaged in unfair competition. |
Publisher Council on Foreign Relations
Release Date Last Updated: May 21, 2013
Scope of the Challenge
Oceans are the source of life on earth. They shape the climate, feed the world, and cleanse the air we breathe. They are vital to our economic well being, ferrying roughly 90 percent of global commerce, housing submarine cables, and providing one-third of traditional hydrocarbon resources (as well as new forms of energy such as wave, wind, and tidal power). But the oceans are increasingly threatened by a dizzying array of dangers, from piracy to climate change. To be good stewards of the oceans, nations around the world need to embrace more effective multilateral governance in the economic, security, and environmental realms.
The world's seas have always been farmed from top to bottom. New technologies, however, are making old practices unsustainable. When commercial trawlers scrape the sea floor, they bulldoze entire ecosystems. Commercial ships keep to the surface but produce carbon-based emissions. And recent developments like offshore drilling and deep seabed mining are helping humans extract resources from unprecedented depths, albeit with questionable environmental impact. And as new transit routes open in the melting Arctic, this once-forgotten pole is emerging as a promising frontier for entrepreneurial businesses and governments.
But oceans are more than just sources of profit—they also serve as settings for transnational crime. Piracy, drug smuggling, and illegal immigration all occur in waters around the world. Even the most sophisticated ports struggle to screen cargo, containers, and crews without creating regulatory friction or choking legitimate commerce. In recent history, the United States has policed the global commons, but growing Indian and Chinese blue-water navies raise new questions about how an established security guarantor should accommodate rising—and increasingly assertive—naval powers.
And the oceans themselves are in danger of environmental catastrophe. They have become the world's garbage dump—if you travel to the heart of the Pacific Ocean, you'll find the North Pacific Gyre, where particles of plastic outweigh plankton six to one. Eighty percent of the world's fish stocks are depleted or on the verge of extinction, and when carbon dioxide is released into the atmosphere, much of it is absorbed by the world's oceans. The water, in response, warms and acidifies, destroying habitats like wetlands and coral reefs. Glacial melting in the polar regions raises global sea levels, which threatens not only marine ecosystems but also humans who live on or near a coast. Meanwhile, port-based megacities dump pollution in the ocean, exacerbating the degradation of the marine environment and the effects of climate change.
Threats to the ocean are inherently transnational, touching the shores of every part of the world. So far, the most comprehensive attempt to govern international waters produced the United Nations Convention on the Law of the Sea (UNCLOS). But U.S. refusal to join the convention, despite widespread bipartisan support, continues to limit its strength, creating a leadership vacuum in the maritime regime. Other states that have joined the treaty often ignore its guidelines or fail to coordinate policies across sovereign jurisdictions. Even if it were perfectly implemented, UNCLOS is now thirty years old and increasingly outdated.
Important initiatives—such as local fishery arrangements and the United Nations Environment Program Regional Seas Program—form a disjointed landscape that lacks legally-binding instruments to legitimize or enforce their work. The recent UN Conference on Sustainable Development ("Rio+20") in Rio de Janeiro, Brazil, convened over one hundred heads of state to assess progress and outline goals for a more sustainable "blue-green economy." However, the opportunity to set actionable targets to improve oceans security and biodiversity produced few concrete outcomes. As threats to the oceans become more pressing, nations around the world need to rally to create and implement an updated form of oceans governance.
Oceans Governance: Strengths and Weaknesses
Overall assessment: A fragmented system
In 1982, the United Nations Convention on the Law of the Sea (UNCLOS) established the fundamental legal principles for ocean governance. This convention, arguably the largest and most complex treaty ever negotiated, entered into force in 1994. Enshrined as a widely accepted corpus of international common law, UNCLOS clearly enumerates the rights, responsibilities, and jurisdictions of states in their use and management of the world's oceans. The treaty defines "exclusive economic zones" (EEZs), which is the coastal water and seabed—extending two hundred nautical miles from shore—over which a state has special rights over the use of marine resources; establishes the limits of a country's "territorial sea," or the sovereign territory of a state that extends twelve nautical miles from shore; and clarifies rules for transit through "international straits." It also addresses—with varying degrees of effectiveness—resource division, maritime traffic, and pollution regulation, as well as serves as the principal forum for dispute resolution on ocean-related issues. To date, 162 countries and the European Union have ratified UNCLOS.
UNCLOS is a remarkable achievement, but its resulting oceans governance regime suffers several serious limitations. First, the world's leading naval power, the United States, is not party to the convention, which presents obvious challenges to its effectiveness—as well as undermines U.S. sovereignty, national interests, and ability to exercise leadership over resource management and dispute resolution. Despite the myriad military, economic, and political benefits offered by UNCLOS, a small but vocal minority in the United States continues to block congressional ratification.
Second, UNCLOS is now thirty years old and, as a result, does not adequately address a number of emerging and increasingly important international issues, such as fishing on the high seas—a classic case of the tragedy of the commons—widespread maritime pollution, and transnational crime committed at sea.
Third, both UNCLOS and subsequent multilateral measures have weak surveillance, capacity-building, and enforcement mechanisms. Although various UN bodies support the instruments created by UNCLOS, they have no direct role in their implementation. Individual states are responsible for ensuring that the convention's rules are enforced, which presents obvious challenges in areas of overlapping or contested sovereignty, or effectively stateless parts of the world. The UN General Assembly plays a role in advancing the oceans agenda at the international level, but its recommendations are weak and further constrained by its lack of enforcement capability.
Organizations that operate in conjunction with UNCLOS—such as the International Maritime Organization (IMO), the International Tribunal on the Law of the Sea (ITLOS), and the International Seabed Authority (ISA)—play an important role to protect the oceans and strengthen oceans governance. The IMO has helped reduce ship pollution to historically low levels, although it can be slow to enact new policy on issues such as invasive species, which are dispersed around the world in ballast water. Furthermore, ITLOS only functions if member states are willing to submit their differences to its judgment, while the ISA labors in relative obscurity and operates under intense pressure from massive commercial entities.
Fourth, coastal states struggle to craft domestic policies that incorporate the many interconnected challenges faced by oceans, from transnational drug smuggling to protecting ravaged fish stocks to establishing proper regulatory measures for offshore oil and gas drilling. UNCLOS forms a solid platform on which to build additional policy architecture, but requires coastal states to first make comprehensive oceans strategy a priority—a goal that has remained elusive thus far.
Fifth, the system is horizontally fragmented and fails to harmonize domestic, regional, and international policies. Domestically, local, state, and federal maritime actors rarely coordinate their agendas and priorities. Among the handful of countries and regional organizations that have comprehensive ocean policies—including Australia, Canada, New Zealand, Japan, the European Union, and most recently the United States—few synchronize their activities with other countries. The international community, however, is attempting to organize the cluttered oceans governance landscape. The UN Environmental Programme Regional Seas Program works to promote interstate cooperation for marine and coastal management, albeit with varying degrees of success and formal codification. Likewise, in 2007 the European Union instituted a regional Integrated Maritime Policy (IMP) that addresses a range of environmental, social, and economic issues related to oceans, as well as promotes surveillance and information sharing. The IMP also works with neighboring partners to create an integrated oceans policy in places such as the Arctic, the Baltic, and the Mediterranean.
Lastly, there is no global evaluation framework to assess progress. No single institution is charged with monitoring and collecting national, regional, and global data on the full range of oceans-related issues, particularly on cross-cutting efforts. Periodic data collecting does take place in specific sectors, such as biodiversity conservation, fisheries issues, and marine pollution, but critical gaps remain. The Global Ocean Observing System is a promising portal for tracking marine and ocean developments, but it is significantly underfunded. Without concrete and reliable data, it is difficult to craft effective policies that address and mitigate emerging threats.
Despite efforts, oceans continue to deteriorate and a global leadership vacuum persists. Much work remains to modernize existing institutions and conventions to respond effectively to emerging threats, as well as to coordinate national actions within and across regions. The June 2012 United Nations Conference on Sustainable Development , also known as Rio+20, identified oceans (or the "blue economy") as one of the seven priority areas for sustainable development. Although experts and activists hoped for a new agreement to strengthen the sustainable management and protection of oceans and address modern maritime challenges such as conflicting sovereignty claims, international trade, and access to resources, Rio+20 produced few concrete results.
Maintaining freedom of the seas: Guaranteed by U.S. power, increasingly contested by emerging states
The United States polices every ocean throughout the world. The U.S. navy is unmatched in its ability to provide strategic stability on, under, and above the world's waters. With almost three hundred active naval ships and almost four thousand aircraft, its battle fleet tonnage is greater than the next thirteen largest navies combined. Despite recently proposed budget cuts to aircraft carriers, U.S. naval power continues to reign supreme.
The United States leverages its naval capabilities to ensure peace, stability, and freedom of access. As Great Britain ensured a Pax Britannicain the nineteenth century, the United States presides over relatively tranquil seas where global commerce is allowed to thrive. In 2007, the U.S. Navy released a strategy report that called for "cooperative relationships with more international partners" to promote "greater collective security, stability, and trust."
The United States pursues this strategy because it has not faced a credible competitor since the end of the Cold War. And, thus far, emerging powers have largely supported the U.S. armada to ensure that the oceans remain open to commerce. However, emerging powers with blue-water aspirations raise questions about how U.S. naval hegemony will accommodate new and assertive fleets in the coming decades. China, for instance, has been steadily building up its naval capabilities over the past decade as part of its "far sea defense" strategy. It unveiled its first aircraft carrier in 2010, and is investing heavily in submarines outfitted with ballistic missiles. At the same time, India has scaled up its military budget by 64 percent since 2001, and plans to spend nearly $45 billion over the next twenty years on its navy.
Even tensions among rising powers could prove problematic. For example, a two-month standoff between China and the Philippines over a disputed region of the South China Sea ended with both parties committing to a "peaceful resolution."China, Taiwan, Vietnam, Malaysia, Brunei, and the Philippines have competing territorial and jurisdictional claims to the South China Sea, particularly over rights to exploit its potentially vast oil and gas reserves. Control over strategic shipping lanes and freedom of navigation are also increasingly contested, especially between the United States and China.
Combating illicit trafficking: Porous, patchy enforcement
In addition to being a highway for legal commerce, oceans facilitate the trafficking of drugs, weapons, and humans, which are often masked by the flow of licit goods. Individual states are responsible for guarding their own coastlines, but often lack the will or capacity to do so. Developing countries, in particular, struggle to coordinate across jurisdictions and interdict. But developed states also face border security challenges. Despite its commitment to interdiction, the United States seizes less than 20 percent of the drugs that enter the country by maritime transport.
The United Nations attempts to combat the trafficking of drugs, weapons, and humans at sea. Through the Container Control Program (PDF), the UN Office on Drugs and Crime (UNODC) assists domestic law enforcement in five developing countries to establish effective container controls to prevent maritime drug smuggling. The UNODC also oversees UN activity on human trafficking, guided by two protocols to the UN Convention on Transnational Organized Crime. Although UN activity provides important groundwork for preventing illicit maritime trafficking, it lacks monitoring and enforcement mechanisms and thus has a limited impact on the flow of illegal cargo into international ports. Greater political will, state capacity, and multilateral coordination will be required to curb illicit maritime trafficking.
New ad hoc multilateral arrangements are a promising model for antitrafficking initiatives. The International Ship and Port Facility Security Code, for instance, provides a uniform set of measures to enhance the security of ships and ports. The code helps member states control their ports and monitor both the people and cargo that travel through them. In addition, the U.S.-led Proliferation Security Initiative facilitates international cooperation to interdict ships on the high seas that may be carrying illicit weapons of mass destruction, ballistic missiles, and related technology. Finally, the Container Security Initiative (CSI), also spearheaded by the United States, attempts to prescreen all containers destined for U.S. ports and identify high-risk cargo (for more information, see section on commercial shipping).
One way to combat illicit trafficking is through enhanced regional arrangements, such as the Paris Memorandum of Understanding on Port State Control. This agreement provides a model for an effective regional inspections regime, examining at least 25 percent of ships that enter members' ports for violations of conventions on maritime safety. Vessels that violate conventions can be detained and repeat offenders can be banned from the memorandum's area. Although the agreement does not permit searching for illegal cargo, it does show how a regional inspections regime could be effective at stemming illegal trafficking.
Securing commercial shipping: Global supply chains at risk
Global shipping is incredibly lucrative, but its sheer scope and breadth presents an array of security and safety challenges. The collective fleet consists of approximately 50,000 ships registered in more than 150 nations. With more than one million employees, this armada transports over eight billion tons (PDF) of goods per year—roughly 90 percent of global trade. And the melting Arctic is opening previously impassable trade routes; in 2009, two German merchant vessels traversed the Northeast Passage successfully for the first time in recent history. But despite impressive innovations in the shipping industry, maritime accidents and attacks on ships still occur frequently, resulting in the loss of billions of dollars of cargo. Ensuring the safety and security of the global shipping fleet is essential to the stability of the world economy.
Internationally, the International Maritime Organization (IMO) provides security guidelines for ships through the Convention on the Safety of Life at Sea, which governs everything from construction to the number of fire extinguishers on board. The IMO also aims to prevent maritime accidents through international standards for navigation and navigation equipment, including satellite communications and locating devices. Although compliance with these conventions has been uneven, regional initiatives such as the Paris Memorandum of Understanding have helped ensure the safety of international shipping.
In addition, numerous IMO conventions govern the safety of container shipping, including the International Convention on Safe Containers, which creates uniform regulations for shipping containers, and the International Convention on Load Lines, which determines the volume of containers a ship can safely hold. However, these conventions do not provide comprehensive security solutions for maritime containers, and illegal cargo could be slipped into shipping containers during transit. Since 1992, the IMO has tried to prevent attacks on commercial shipping through the Convention for the Suppression of Unlawful Acts against the Safety of Maritime Navigation, which provides a legal framework for interdicting, detaining, and prosecuting terrorists, pirates, and other criminals on the high seas.
In reality, most enforcement efforts since the 9/11 attacks have focused on securing ports to prevent the use of a ship to attack, rather than to prevent attacks on the ships themselves. Reflecting this imperative, the IMO, with U.S. leadership, implemented the International Ship and Port Facility Security Code (ISPS) in 2004. This code helped set international standards for ship security, requiring ships to have security plans and officers. However, as with port security, the code is not obligatory and no clear process to audit or certify ISPS compliance has been established. Overall, a comprehensive regime for overseeing the safety of international shipping has not been created.
The United States attempts to address this vulnerability through the Container Security Initiative (CSI), which aims to prescreen all containers destined for the United States, and to isolate those that pose a high-security risk before they are in transit. The initiative, which operates in fifty-eight foreign ports, covers more than 86 percent of container cargo en route to the United States. Several international partners and organizations, including the European Union, the Group of Eight, and the World Customs Organization, have expressed interest in modeling security measures for containerized cargo based on the CSI model. Despite these efforts, experts estimate that only 2 percent of containers destined for U.S. ports are actually inspected.
Confronting piracy: Resurgent scourge, collective response
After the number of attacks reached a record high in 2011, incidences of piracy dropped 28 percent in the first three months of 2012. Overall, the number of worldwide attacks decreased from 142 to 102 cases, primarily due to international mobilization and enhanced naval patrols off the coast of Somalia. However, attacks intensified near Nigeria and Indonesia as pirates shifted routes in response to increased policing, raising fresh concerns over the shifting and expanding threat of piracy. In addition to the human toll, piracy has significant economic ramifications. According to a report by the nonprofit organization Oceans Beyond Piracy, Somali piracy cost the global economy nearly $7 billion in 2011. Sustained international coordination and cooperation is essential to preventing and prosecuting piracy.
Recognizing this imperative, countries from around the world have shown unprecedented cooperation to combat piracy, particularly near the Gulf of Aden. In August 2009, the North Atlantic Treaty Organization commenced Operation Ocean Shield in the horn of Africa, where piracy increased close to 200 percent between 2007 and 2009. This effort built upon Operation Allied Protector and consisted of two standing maritime groups with contributions from allied nations. Although the efforts concentrate on protecting ships passing through the Gulf of Aden, they also renewed focus on helping countries, specifically Somalia, prevent piracy and secure their ports. Meanwhile, the United States helped establish Combined Task Force 151 to coordinate the various maritime patrols in East Africa. Other countries including Russia, India, China, Saudi Arabia, Malaysia, and South Korea, have also sent naval vessels to the region.
At the same time, regional organizations have also stepped up antipiracy efforts. The Regional Cooperation Agreement on Combating Piracy and Armed Robbery against Ships in Asia was the first such initiatives, and has been largely successful in facilitating information-sharing, cooperation between governments, and interdiction efforts. And in May 2012, the European Union naval force launched its first air attack against Somali pirates' land bases, the first strike of its kind by outside actors to date.
Like individual countries, international institutions have condemned piracy and legitimized the use of force against pirates. In June 2008, the UN Security Council unanimously passed Resolution 1816, encouraging greater cooperation in deterring piracy and asking countries to provide assistance to Somalia to help ensure coastal security. This was followed by Resolution 1846, which allowed states to use "all necessary means" to fight piracy off the coast of Somalia. In Resolution 1851, the UN Security Council legitimized the use of force on land as well as at sea to the same end. Outside the UN, watchdogs such as the International Maritime Bureau, which collects information on pirate attacks and provides reports on the safety of shipping routes, have proven successful in increasing awareness, disseminating information, and facilitating antipiracy cooperation.
However, such cooperative efforts face several legal challenges. The United States has not ratified the UN Convention on the Law of the Sea (UNCLOS), which governs crimes, including piracy, in international waters. More broadly, the international legal regime continues to rely on individual countries to prosecute pirates, and governments have been reluctant to take on this burden. Accordingly, many pirates are apprehended, only to be quickly released. In addition, many large commercial vessels rely on private armed guards to prevent pirate attacks, but the legal foundations governing such a force are shaky at best.
National governments have redoubled efforts to bring pirates to justice as well. In 2010, the United States held its first piracy trial since its civil war, soon followed by Germany's first trial in over four hundred years. Other agreements have been established to try pirates in nearby countries like Kenya, such as the UNODC Trust Fund to Support the Initiatives of States to Counter Piracy of the Coast of Somalia, established in January 2010. Under the mandate of the Contact Group on Piracy off the Coast of Somalia, the fund aims to defray the financial capital required from countries like Kenya, Seychelles, and Somalia to prosecute pirates, as well as to increase awareness within Somali society of the risk associated with piracy and criminal activity. Future efforts to combat piracy should continue to focus on enhancing regional cooperation and agreements, strengthening the international and domestic legal instruments necessary to prosecute pirates, and addressing the root causes of piracy.
Reducing marine pollution and climate change: Mixed progress
Pollution has degraded environments and ravaged biodiversity in every ocean. Much contamination stems from land-based pollutants, particularly along heavily developed coastal areas. The UN Environment Program (UNEP) Regional Seas Program has sponsored several initiatives to control pollution, modeled on a relatively successful program in the Mediterranean Sea. In 1995, states established the Global Program of Action for the Protection of the Marine Environment from Land-Based Activities, which identifies sources of land-based pollution and helps states establish priorities for action. It has been successful in raising awareness about land-based pollution and offering technical assistance to regional implementing bodies, which are so often starved for resources. More recently, 193 UN member states approved the Nagoya Protocol on biodiversity, which aims to halve the marine extinction rate by 2020 and extend protection to 10 percent of the world's oceans.
Shipping vessels are also a major source of marine pollution. Shipping is the most environmentally friendly way to transport bulk cargoes, but regulating maritime pollution remains complicated because of its inherently transnational nature. Shipping is generally governed by the International Maritime Organization (IMO), which regulates maritime pollution through the International Convention for the Prevention of Pollution from Ships (MARPOL). States are responsible for implementing and enforcing MARPOL among their own fleets to curb the most pernicious forms of maritime pollution, including oil spills, particulate matter such as sulfur oxide (SOx) and nitrous oxide (NOx), and greenhouse gas emissions. Port cities bear the brunt of air pollution, which devastates local air quality because most ships burn bunker fuel (the dirtiest form of crude oil). The IMO's Marine Environmental Protection Committee has also taken important steps to reduce SOx and NOx emissions by amending the MARPOL guidelines to reduce particulate matter from ships. Despite such efforts, a 2010 study (PDF) from the Organization for Economic Development and Cooperation found that international shipping still accounts for nearly 3 percent of all greenhouse gasses.
The IMO has achieved noteworthy success in reducing oil spilled into the marine environment. Despite a global shipping boom, oil spills are at an all-time low. The achievements of the IMO have been further strengthened by commitments by the Group of Eight to cooperate on oil pollution through an action plan that specifically targets pollution prevention for tankers. The IMO should strive to replicate this success in its efforts to reduce shipping emissions.
Climate change is also exacerbating environmental damage. In June 2009, global oceans reached their highest recordedaverage temperature: 17 degrees Celsius. As the world warms, oceans absorb increased levels of carbon dioxide, which acidifies the water and destroys wetlands, mangroves, and coral reefs—ecosystems that support millions of species of plants and animals. According to recent studies, ocean acidity could increase by more than 150 percent by 2050 if counteracting measures are not taken immediately. Moreover, melting ice raises sea levels, eroding beaches, flooding communities, and increasing the salinity of freshwater bodies. And the tiny island nation of the Maldives, the lowest country in the world, could be completely flooded if sea levels continue to rise at the same rate.
Individual states are responsible for managing changes in their own marine climates, but multilateral efforts to mitigate the effect of climate change on the oceans have picked up pace. In particular, the UNEP Regional Seas Program encourages countries sharing common bodies of water to coordinate and implement sound environmental policies, and promotes a regional approach to address climate change.
Sustainable fisheries policies on the high seas: An ecological disaster
States have the legal right to regulate fishing in their exclusive economic zones (EEZs), which extend two hundred nautical miles from shore—and sometimes beyond, in the case of extended continental shelves. But outside the EEZs are the high seas, which do not fall under any one country's jurisdiction. Freedom of the high seas is critical to the free flow of global commerce, but spells disaster for international fisheries in a textbook case of the tragedy of the commons. For years, large-scale fishing vessels harvested fish as fast as possible with little regard for the environmental costs, destroying 90 percent of the ocean's biomass in less than a century. Overall, fisheries suffer from two sets of challenges: ineffective enforcement capacity and lack of market-based governance solutions to remedy perverse incentives to overfish.
Although there are numerous international and multilateral mechanisms for fisheries management, the system is marred by critical gaps and weaknesses exploited by illegal fishing vessels. Articles 117 and 118 of the UN Convention on the Law of the Sea (UNCLOS) enumerate the specific fisheries responsibilities of state parties, placing the onus on national governments to form policies and regional agreements that ensure responsible management and conservation of fish stocks in their respective areas. UNCLOS was further strengthened by the UN Fish Stocks Agreement (FSA), which called for a precautionary approach toward highly migratory and straddling fish stocks that move freely in and out of the high seas. Seventy-eight countries have joined the FSA thus far, and a review conference in May 2010 was hailed as a success due to the passage of Port State Measures (PSMs) to combat illegal, unreported, and unregulated (IUU) fishing. Yet fish stocks have continued to stagnate or decline to dangerously low levels, and the PSMs have largely failed to prevent IUU operations.
Regional fishery bodies (RFBs) are charged with implementation and monitoring. The RFBs provide guidelines and advice on a variety of issues related to fishing, including total allowable catch, by-catch, vessel monitoring systems, areas or seasons closed for fishing, and recording and reporting fishery statistics. However, only a portion of these bodies oversee the management of their recommendations, and some RFBs allow members to unilaterally dismiss unfavorable decisions. Additionally, RFBs are not comprehensive in their membership and, for the most part, their rules do not apply to vessels belonging to a state outside the body.
Even when regional bodies make a binding decision on a high-seas case, implementation hinges on state will and capacity. In 2003, the UN General Assembly established a fund to assist developing countries with their obligations to implement the Fish Stocks Agreement through RFBs. The overall value of the fund remains small, however, and countries' compliance is often constrained by resource scarcity. This results in spotty enforcement, which allows vessels to violate international standards with impunity, particularly off the coasts of weak states. Migratory species like blue fin tuna are especially vulnerable because they are not confined by jurisdictional boundaries and have high commercial value.
Some of the RFBs with management oversight, such as the Commission for the Conservation of Antarctic Marine Living Resources and the South East Atlantic Fisheries Organization, have been relatively effective in curbing overfishing. They have developed oversight systems and specific measures to target deep-water trawl fishing and illegal, unreported, and unregulated fishing in the high seas. Many regional cooperative arrangements, however, continue to suffer from weak regulatory authority. At the same time, some regions like the central and southwest Atlantic Ocean lack RFBs. Some have suggested filling the void with market-based solutions like catch shares, which could theoretically alter the incentives toward stewardship. Catch shares (also known as limited access privilege programs) reward innovation and help fisheries maximize efficiency by dedicating a stock of fish to an individual fisherman, community, fishery association, or an individual state. Each year before the beginning of fishing season, commercial fishermen would know how much fish they are allowed to catch. They would then be allowed to buy and sell shares to maximize profit. By incorporating free-market principles, fisheries could reach a natural equilibrium at a sustainable level. According to research, more sustainable catch shares policies could increase the value of the fishing industry by more than $36 billion. Although allocating the shares at the domestic—much less international—level remains problematic, the idea reflects of the kind of policy work required to better manage the global commons.
Managing the Arctic: At a crossroads
Arctic ice is melting at unprecedented rates. At this pace, experts estimate that the Arctic could be seasonally ice free by 2040, and possibly much earlier. As the ice recedes and exposes valuable new resources, multilateral coordination will become even more important among states (and indigenous groups) jockeying for position in the region.
The melting ice is opening up potentially lucrative new sea routes and stores of natural resources. Since September 2009, cargo ships have been able to traverse the fabled Northwest and Northeast Passages, which are significantly shorter than traditional routes around the capes or through the canals. Widening sea routes also means that fishing fleets can travel north in search of virgin fishing stock, and that cruise ships can carry tourists chasing a last glimpse of the disappearing ice. At the same time, untapped resources such as oil, natural gas, rare earth minerals, and massive renewable wind, tidal, and geothermal energy hold enormous potential. In a preliminary estimate, the U.S. Geographic Society said that the Arctic could hold 22 percent of the world's hydrocarbon resources, including 90 billion barrels of oil and 1,670 trillion cubic feet of natural gas. Beyond oil and gas, the Arctic has valuable mineral commodities such as zinc, nickel, and coal.
But new opportunities in the Arctic also portend new competition among states. In August 2007, Russia symbolically planted a flag on the Arctic floor, staking a claim to large chunks of Arctic land. Other Arctic powers including the United States, Canada, Norway, and Denmark have also laid geographical claims. The European Union crafted a new Arctic policy, and China sent an icebreaker on three separate Arctic expeditions. Each country stands poised to grab new treasure in this increasingly important geostrategic region.
The UN Convention on the Law of the Sea (UNCLOS) is a solid foundation on which to build and coordinate national Arctic policies, especially articles 76 and 234, which govern the limits of the outer continental shelf (OCS) and regulate activities in ice-covered waters, respectively. However, there remains a formidable list of nagging sovereignty disputes that will require creative bilateral and multilateral resolutions. The Arctic Council, a multilateral forum comprising eight Arctic nations, has recently grown in international prominence, signing a legally binding treaty on search and rescue missions in May 2011 and drawing high-level policymakers to its meetings. While these are significant first steps, the forum has yet to address other issues such as overlapping OCS claims, contested maritime boundaries, and the legal status of the Northwest Passage and the Northern Sea Route.
U.S. Ocean Governance Issues
The United States championed many of the most important international maritime organizations over the past fifty years. It helped shape the decades-long process of negotiating the United Nations Convention on the Law of the Sea (UNCLOS) and has played a leading role in many UNCLOS-related bodies, including the International Maritime Organization. It has also served as a driving force behind regional fisheries organizations and Coast Guard forums. Domestically, the United States has intermittently been at the vanguard of ocean policy, such as the 1969 Stratton Commission report, multiple conservation acts in the 1970s, the Joint Ocean Commission Initiative, and, most recently, catch limits on all federally-managed fish species. The U.S.-based Woods Hole Oceanographic Institution and the Monterrey Bay Research Institute have long been leaders in marine science worldwide. And from a geopolitical perspective, the U.S. Navy secures the world's oceans and fosters an environment where global commerce can thrive.
Yet the United States lags behind on important issues, most notably regarding its reluctance to ratify UNCLOS. And until recently, the United States did not have a coherent national oceans policy. To address this gap, U.S. president Barack Obama created the Ocean Policy Task Force in 2009 to coordinate maritime issues across local, state, and federal levels, and to provide a strategic vision for how oceans should be managed in the United States. The task force led to the creation of a National Ocean Council, which is responsible for "developing strategic action plans to achieve nine priority objectives that address some of the most pressing challenges facing the ocean, our coasts, and Great Lakes." Although it has yet to make serious gains (PDF), this comprehensive oceans policy framework could help clear the way for the spadework of coordinating U.S. ocean governance and harmonizing international efforts.
Should the United States ratify the UN Convention on the Law of the Sea?
Yes: The UN Convention on the Law of the Sea (UNCLOS), which created the governance framework that manages nearly three-quarters of the earth's surface, has been signed and ratified by 162 countries and the European Union. But the United States remains among only a handful of countries to have signed but not yet ratified the treaty—even though it already treats many of the provisions as customary international law. Leaders on both sides of the political aisle as well as environmental, conservation, business, industry, and security groups have endorsed ratification in order to preserve national security interests and reap its myriad benefits, such as securing rights for U.S. commercial and naval ships and boosting the competitiveness of U.S. companies in seafaring activities. Notably, all of the uniformed services—and especially the U.S. Navy—strongly support UNCLOS because its provisions would only serve to strengthen U.S. military efforts. By remaining a nonparty, the United States lacks the credibility to promote its own interests in critical decision-making forums as well as bring complaints to an international dispute resolution body.
No: Opponents argue that ratifying the treaty would cede sovereignty to an ineffective United Nations and constrain U.S. military and commercial activities. In particular, critics object to specific provisions including taxes on activities on outer continental shelves; binding dispute settlements; judicial activism by the Law of the Sea Tribunal, especially with regard to land-based sources of pollution; and the perceived ability of UNCLOS to curtail U.S. intelligence-gathering activities. Lastly, critics argue that because UNCLOS is already treated as customary international law, the United States has little to gain from formal accession.
Should the United States lead an initiative to expand the Container Security Initiative globally?
Yes: Some experts say the only way to secure a global economic system is to implement a global security solution. The U.S.-led Container Security Initiative (CSI) helps ensure that high-risk containers are identified and isolated before they reach their destination. Fifty-eight countries are already on board with the initiative, and many others have expressed interest in modeling their own security measures on the CSI. The World Customs Organization called on its members to develop programs based on the CSI, and the European Union agreed to expand the initiative across its territory. With its robust operational experience, the United States is well positioned to provide the technical expertise to ensure the integrity of the container system.
No: Opponents maintain that the United States can hardly commit its tax dollars abroad for a global security system when it has failed to secure its own imports. To date, more than $800 million and considerable diplomatic energy has been invested in CSI to expand the program to fifty-eight international ports, where agents are stationed to screen high-risk containers. Given the scale of world trade, the United States imports more than 10 million containers annually, and only a handful of high-risk boxes can be targeted for inspection. After huge expenditures and years of hard work to expand this program after September 11, 2001, only about 86 percent of the cargo that enters the United States transits through foreign ports covered under CSI, and of that, only about 1 percent is actually inspected (at a cost to the U.S. taxpayer of more than $1,000 per container). Despite congressional mandates to screen all incoming containers, critics say that costs make implementing this mandate virtually impossible. The limited resources the United States has available, they argue, should be invested in protecting imports bound specifically for its shores.
Should the United States be doing more to address the drastic decline in the world's fisheries?
Yes: Advocates say that the further demise of global fish stocks, beyond being a moral burden, undermines the commercial and national security interests of the United States. Depleting fish stocks are driven in large part by the prevalence of illegal, unreported, and unregulated (IUU) fishing and the overcapitalization of the global commercial fishing fleet from domestic subsidies. To protect domestic commercial fisheries and the competitiveness of U.S. exports in the international seafood market, the United States should enhance efforts by the National Oceanic and Atmospheric Administration to manage, enforce, and coordinate technical assistance for nations engaging in IUU fishing.
Domestically, the United States has taken important steps to address the critical gaps in fisheries management. In 2012, it became the first country to impose catch limits on all federally-managed fish species. Some species like the mahi mahi will be restricted for the first time in history. Many environmental experts hailed the move as a potential model for broader regional and international sustainable fisheries policy. To capitalize on such gains, the United States should aggressively work to reduce fishing subsidies in areas such as Europe that promote overcapitalization and thus global depletion of fish stocks. The United States could also promote market-based mechanisms, like catch shares and limited access privilege programs, to help fishermen and their communities curb overfishing and raise the value of global fisheries by up to $36 billion.
No: Critics argue that fisheries management is by and large a domestic issue, and that the United States has little right to tell other nations how to manage their own resources, particularly when such measures could harm local economies. They contend that the science behind overfishing is exaggerated, as are the warnings about the consequences of an anticipated fisheries collapse. Existing conventions like the 1995 Fish Stock Agreement already go far enough in addressing this issue. Any additional efforts, they contend, would be a diplomatic overreach, as well as an excessive burden on a struggling commercial fishing industry. Critics also question how market-based mechanisms, such as catch-shares, would be distributed, traded, and enforced, warning that they would lead to speculative bubbles.
Should the United States push for a more defined multilateral strategy to cope with the melting Arctic?
Yes: The melting Arctic holds important untapped political, strategic, and economic potential for the U.S. government, military, and businesses. This emerging frontier could potentially support a variety of economic activities, including energy exploration, marine commerce, and sustainable development of new fisheries. Countries such as Russia, Canada, Norway, and China have already made claims to the region, yet the United States remains on the sideline without a comprehensive Arctic strategy. The UN Convention on the Law of the Sea (UNCLOS) remains the premier forum of negotiating and arbitrating disputes over contested territory. As a nonparty, however, the United States loses invaluable leverage and position. In addition, the U.S. military does not have a single icebreaker, whereas Russia operates over thirty. Experts argue that the U.S. government should also adopt the recently proposed Polar Code, which is a voluntary agreement that "sets structural classifications and standards for ships operating in the Arctic as well as specific navigation and emergency training for those operating in or around ice-covered waters."
No: Opponents argue that Arctic Council activities and the 2009 National Security Presidential Directive, which updated U.S. Arctic polices, are sufficient. Any collaboration with Canada to resolve disputes over the Northwest Passage might undermine freedom of navigation for U.S. naval assets elsewhere, especially in the Strait of Hormuz and the Taiwan Straits, and this national security concern trumps any advantages from collaborating on security, economic, or environmental issues in the Arctic. Last, given the dominant Russian and Canadian Arctic coastlines, future Arctic diplomacy might best be handled bilaterally rather than through broader multilateral initiatives.
April 2013: Japan included in Trans-Pacific Partnership negotiations
Japan agreed to join negotiations over the Trans-Pacific Partnership (TPP), an ambitious free trade agreement between counties along the Pacific rim. Since the broad outline of the agreement was introduced in November 2011, sixteen rounds of negotiations have thus far brought eleven countries together to discuss the TPP. The addition of Japan, a major economic force in the region, as the twelfth participant comes as an important step in creating a robust agreement. Already, the South China Sea is the second-busiest shipping lane in the world, and should the TTP become a reality, transpacific shipping would dramatically increase. The seventeenth round of negotiations will take place in May and the current goal for agreement is October 2013.
March 2013: IMO pledges to support implementation of new code of conduct on piracy
At a ministerial meeting in Cotonou, Benin, the International Maritime Organization (IMO) pledged to support the implementation of a new code of conduct on piracy and other illicit maritime activity. The Gulf of Guinea Code of Conduct, drafted by the Economic Community of Central African States and the Economic Community of West African States, in partnership with the IMO, contains provisions for interdicting sea- and land-based vehicles engaged in illegal activities at sea, prosecuting suspected criminals, and sharing information between state parties. The code builds on several existing frameworks to create a sub-regional coast guard. The agreement is set to open for signature in May 2013.
March 2013: New fishing restrictions on sharks and rays
Delegates attending the annual meeting on the Convention on International Trade in Endangered Species of Wild Fauna and Flora (CITES) voted to place robust export restrictions on five species of sharks and two species of manta rays. Over the past fifty years, the three shark species—the oceanic whitetip, hammerhead, and porbeagle—have declined by more than 70 percent. Although experts cautioned that the new rules would be difficult to enforce in practice, the decision marked an important victory over economic interests, particularly of China and Japan.
January 2013: Philippines to challenge China's maritime claims in South China Sea
The Philippine government announced its intention to take China to an international arbitration tribunal based on claims that China violated the UN Convention on the Law of the Sea. The dispute dates back to mid-2012, when tensions flared over the Scarborough shoal, which is claimed by both countries.
China, Taiwan, Vietnam, Malaysia, Brunei, and the Philippines have competing territorial and jurisdictional claims to the South China Sea, particularly over rights to exploit its potentially vast oil and gas reserves. Control over strategic shipping lanes and freedom of navigation are also increasingly contested, especially between the United States and China.
September 2012: Arctic ice reaches record low
In September 2012, ice in the Arctic Ocean reached an all-time low of 24 percent, shattering the previous record of 29 percent from 2007. The finding not only has implications for climate change and environmental stability, but also for heightened competition among states jockeying for access to critical resources in the region. For the first time, the melting Arctic has exposed troves of natural resources including oil, gas, and minerals, as well as newly accessible shipping routes. The United States, Russia, and several European states already control parts of the Arctic, and China is also an increasing presence.
September 2012: Tensions flare in the East China Sea
In September 2012, Japan purchased three islands in the East China Sea that form part of the Senkaku Islands, known as the Diaoyu Islands to the Chinese. The islands, claimed by both countries, have been controlled by Japan since 1895, but sovereignty remains hotly contested. Following Japan's announcement, protests broke out across China, and Chinese leaders accused Japan of "severely infringing" upon their sovereignty.
In a move to affirm its claim to the islands, China announced its intention to submit their objections to the Commission on the Limits of the Continental Shelf under the UNCLOS, and dispatched patrol ships to monitor the islands. In December 2012, tensions flared after a Chinese small aircraft flew into airspace over the islands, and both countries sent naval vessels to patrol nearby waters. Both sides remain adamant that there is no room for negotiations over their control of the islands, which are in close proximity to strategic shipping routes, fishing grounds, and potentially lucrative oil reserves.
Options for Strengthening Global Ocean Governance
There are a series of measures, both formal and informal, that can be taken to strengthen U.S. and global ocean governance. First, the United States must begin by finally ratifying the UN Convention on the Law of the Sea. On this foundation, the United States should then tap hitherto underused regimes, update twentieth-century agreements to reflect modern ocean challenges, and, in some cases, serve as the diplomatic lead in pioneering new institutions and regimes. These recommendations reflect the views of Stewart M. Patrick, senior fellow and director of the International Institutions and Global Governance Program, and Scott G. Borgerson, former visiting fellow for ocean governance.
In the near term, the United States and its international partners should consider the following steps:
- Ratify UNCLOS
The United States should finally join the UN Convention on the Law of the Sea (UNCLOS), an action that would give it further credibility and make the United States a full partner in global ocean governance. This carefully negotiated agreement has been signed and ratified by 162 countries and the European Union. Yet despite playing a central role shaping UNCLOS's content, the United States has conspicuously failed to join. It remains among only a handful of countries with a coastline, including Syria, North Korea, and Iran, not to have done so.
Emerging issues such as the melting Arctic lend increased urgency to U.S. ratification. By rejecting UNCLOS, the United States is freezing itself out of important international policymaking bodies, forfeiting a seat at decision-making forums critical to economic growth and national security interests. One important forum where the United States has no say is the commission vested with the authority to validate countries' claims to extend their exclusive economic zones, a process that is arguably the last great partitioning of sovereign space on earth. As a nonparty to the treaty, the United States is forgoing an opportunity to extend its national jurisdiction over a vast ocean area on its Arctic, Atlantic, and Gulf coasts—equal to almost half the size of the Louisiana Purchase—and abdicating an opportunity to have a say in deliberations over other nations' claims elsewhere.
Furthermore, the convention allows for an expansion of U.S. sovereignty by extending U.S. sea borders, guaranteeing the freedom of ship and air traffic, and enhancing the legal tools available to combat piracy and illicit trafficking. Potential participants in U.S.-organized flotillas and coalitions rightly question why they should assist the United States in enforcing the rule of law when the United States refuses to recognize the convention that guides the actions of virtually every other nation.
- Coordinate national ocean policies for coastal states
The creation of a comprehensive and integrated U.S. oceans policy should be immediately followed by similar efforts in developing maritime countries, namely Brazil, Russia, India, and China (BRIC) . These so-called BRIC nations will be critical players in crafting domestic ocean policies that together form a coherent tapestry of global governance. Ideally, such emerging powers would designate a senior government official, and in some cases the head of state, to liaison with other coastal states and regional bodies to coordinate ocean governance policies and respond to new threats. Consistent with the Regional Seas Program, the ripest opportunity for these efforts is at the regional level. With UN assistance, successful regional initiatives could then be harmonized and expanded globally.
- Place a moratorium on critically endangered commercial fisheries
Commercial fishing, a multi-billion dollar industry in the United States, is in grave danger. The oceans have been overfished, and it is feared that many fish stocks may not rebound. In the last fifty years, fish that were previously considered inexhaustible have been reduced to alarmingly low levels. Up to 90 percent of large predatory fish are now gone. Nearly half of fish stocks in the world have been fully exploited and roughly one-third have been overexploited. The recent imposition of catch limits on all federally-managed fish species is an important and long overdue first step, which should be expanded and strengthened to a moratorium on the most endangered commercial fisheries, such as the Atlantic blue fin tuna. But tuna is hardly alone in this predicament, and numerous other species are facing the same fate. Policymakers should stand up to intense political pressure and place fishing moratoriums on the most threatened fisheries to give them a chance to rebound. Doing so would be a courageous act that would help rescue collapsing fish while creating a commercially sustainable resource.
In the longer term, the United States and its international partners should consider the following steps:
- Strengthen and update UNCLOS
The UN Convention on the Law of the Sea (UNCLOS) and related agreements serve as the bedrock of international ocean policy. However, UNCLOS is thirty years old. If it is to remain relevant and effective, it must be strengthened and updated to respond to emerging threats such as transnational crime and marine pollution, as well as employing market-based principles of catch shares to commercial fisheries, especially in the high seas. Lastly, UNCLOS Article 234, which applies to ice-covered areas, should be expanded to better manage the opening Arctic, which will be an area of increasing focus and international tension over the coming years.
The international community should also counter the pressure of coastal states that unilaterally seek to push maritime borders seaward, as illustrated by China's claim to all of the South China Sea. Additionally, states should focus on using UNCLOS mechanisms to resolve nagging maritime conflicts, such as overlapping exclusive economic zones from extended continental shelf claims, and sovereignty disputes, such as that of the Spratly and Hans Islands.
- Bolster enforcement capacity
Many ocean-related governance issues have shortcomings not because rules for better management do not exist, but because weak states cannot enforce them. A failure in the oversight of sovereign waters inevitably leads to environmental degradation and, in cases like Somalia, can morph into problems with global implications, such as piracy. Accordingly, the international community should help less developed coastal states build the capacity to enforce (1) fisheries rules fleets; (2) International Convention for the Prevention of Pollution From Ships regulations to reduce ocean dumping and pollution; (3) other shipping regulations in states with open registries such as Liberia, Panama, Malta, and the Marshall Islands; (4) and existing mandates created to stop illicit trafficking. Developed countries should also help less developed areas monitor environmental variables such as acidification, coral reefs, and fisheries. |
On the Lap of the Mighty Sagarmatha - Solu Khumbu or Everest region
The major mountains are the Mt.Everest, Mt.Lhotse, Cho Oyu, Nuptse, Pumori, Ama Dablam, Thamserku, Kantega, Mera Peak and Island Peak.
Mt. Everest, which is part of the Himalaya range, is located on the border between Nepal and Tibet. Rising to a height of 8848m, the world’s highest mountain was named in 1865 after Sir George Everest. The mountain got its Nepali name Sagarmatha during the 1960s, when the Government of Nepal gave the mountain the official Nepali name. In sanskrit Sagarmatha means "mother of the universe”. The Tibetan name for Mount Everest is Chomolungma or Qomolangma, which means “Goddess Mother of the Snows". Climbers wishing to scale the peak have to obtain an expensive permit from the Nepal Government, often costing more than $25,000 (USD) per person. Base Camp, which serves as a resting area and base of operations for climbers organizing their attempts for the summit, is located on the Khumbu glacier at an elevation of 5300 m (17,400 ft); it receives an average of 450 mm (18 in) of precipitation a year. The climate of Mount Everest is extreme In July, the warmest month, the average summit temperature is -19° C (-2° F).
When George Mallory, the British climber was asked why he wanted to climb Everest he replied ‘Because it is there’. After two unsuccessful attempts, in 1924 he again tried to climb the peak with Andrew Irvine. They started on June 8, 1924 to scale the summit via the north col route and never returned. Their bodies were later discovered by the Mallory and Irvine Research Expedition near the old Chinese camp in 1999. Edmund Hillary, a New Zealander and Sherpa Tenzing Norgay from Nepal were the first two climbers to set foot on the summit of Mt.Everest. They reached the summit at 11:30 a.m. on May 29, 1953 by climbing through the South Col Route. More than 300 climbers have scaled the highest mountain since then. Also there have been more than 100 deaths on the mountain where conditions are so difficult that most corpses have been left where they fell, some of them visible from standard climbing routes.
Mt. Lhotse (8516m) is the fourth highest mountain in the world. It lies south of Mt. Everest. It was first climbed by two Swiss climbers F. Luchsinger and E. Reiss in 1956 from the West face. The Czech scaled it via the South face in 1984. An impressive ring of three peaks makes up the Lhotse massif: Lhotse East or Middle, Lhotse and Lhotse Shar. The South Face of Lhotse is one of the largest mountain faces in the world.
Cho Oyu, (8201m) the sixth highest mountain in the world, has gained popularity among climbers just recently. The mountain sits on both sides of the border of Nepal and Tibet, about 30 km. west of Mount Everest. Cho Oyu in Tibetan means "the turquoise goddess ." The south face of Cho Oyu, facing Nepal, is quite steep and difficult, and is rarely climbed. The north side, accessed from Tibet, is more moderate, and there is a relatively safe route to the summit. In the autumn of 1954, an Austrian team made the first ascent via this route.
Ama Dablam (6856m) which means ‘mother’s jewellery box’, in sherpa language is considered to be one of the most beautiful mountains in the world. Seen from below, the mountain looks like a woman with outstretched arms or a woman wearing a long necklace. Ama Dablam lies alongside Everest in the heart of the Khumbu valley. Mt Lhotse, Mt. Makalu, Mt. Cho Oyu and Mt. Everest can be seen at close quarters from Ama dablam.
Nuptse (7,855m.) lies south west of Mt Everest. It is situated in the Khumbu Himal. From the Thyangboche Monastery Nuptse appears as a massive wall guarding the approach to Everest. The name Nup-tse in Tibetan means west-peak. The main ridge, which is separated from Lhotse by a 7556m high saddle, is crowned by seven peaks and goes west-northwest until its steep west-face drops down more than 2300m to the Khumbu-glacier. Nuptse I was first summited by a British expedition on May 16, 1961
Pumori peak at 7145m is just 8 km away from the world’s highest peak Mt.Everest. The ascent to this peak is described as a classic climb in the 7000m peak category. In Tibetan, ‘Pumo’ means girl and ‘Ri’, mountain. The peak was named by George Mallory, the famous English climber who lost his life trying to ascend Everest in 1924. The German climber Gerhard Lenser was the first to reach the summit of Pumori peak in 1962. Pumori is a popular climbing peak and the easiest. The best season to climb this peak is during autumn and spring.
Mera Peak (6,475m) is the highest of Nepal's trekking peaks. By its standard route, it is also the highest peak in Nepal that can be climbed without prior mountaineering experience. It was first climbed on 20 May 1953, by J.O.M. Roberts and Sen Tenzing, from the standard route at Mera La. The mountain lies to the south of Everest, dominating the watershed between the wild and beautiful valleys of the Hinku and Hongu.
Island Peak also known as Imja Tse at 6160m was named by Erick Shipton's group in 1953.It was so named as the peak resembles an island in a sea of ice when observed from Dingboche. The peak was first climbed in 1953 by a British group as preparation for climbing Mt. Everest. Among them one of the climbers was Mr. Tenzing Norgay. The peak is part of the south ridge of Lhotse Shar and the main land forms a semicircle of cliffs that rise to the north of the summits of Nuptse, Lhotse, Middle Peak and Lhotse Shar. Cho Oyu and Makalu lie to the east of the Island Peak. Baruntse, Amphu and Ama Dablam lie to the south.
Lobuche(6,119m) is known as Lhauche among the Locals. It rises above the town of Lhauche which is just a few kilometer from Mt. Everest. The first ascent on this peak was done by Laurice Nielson and Ang Gyalzen Sherpa on 25 April 1984.
Kalapattar is a small mountain 5,545 m (18,500 ft) high on the southern flank of Pumori (7,145 m). It is a trekking peak and every year tourists climb this peak to enjoy the fantastic panoramic views it offers of the Khumbu glacier, the Everest and nearby peaks like Lhotse and Nuptse. To the east, Makalu, Amadablam, Pumori, and Cho Oyu are visible.
Climate, Flora & Fauna
The climate in the Everest region can be divided into four climate zones owing to the gradual rise in altitude. The climatic zones include a forested lower zone, a zone of alpine scrub, the upper alpine zone which includes upper limit of vegetation growth, and the Arctic zone where no plants can grow. The types of plants and animals that are found depend on the altitude. In the lower forested zone, birch, juniper, blue pines, firs, bamboo and rhododendron grow. All vegetation that is found above this zone is shrubs. As the altitude increases, plant life is restricted to lichens and mosses. At an elevation of 5,750m begins the permanent snow line in the Himalayas. From this point there is no sign of greenery or vegetation. A common animal sighted in the higher reaches is the hairy animal yak. Dzopkyo a sterile male crossbreed between a yak and a cow is used to move goods along the trail. Red panda, snow leopard, musk deer, wild yak, and Himalayan black bear are some of the more exotic animals that are found in this region. A variety of birds can be sighted in the lower regions.
Sagarmata (Mt. Everest) National Park
The Sagarmatha National Park is the highest national park in the world. It was formally opened to public in July 19, 1976. The park covers an area of 1,148 sq km. It rises from its lowest point of 2,845 m (9,335 ft) at Jorsale to 8,850 m (29,035 ft) up to the summit of Everest. The park’s area is very rugged and steep, with its terrain cut by deep rivers and glaciers. It includes three peaks higher than 8,000 m, including Mt Everest. In 1979 the park was inscribed as a Natural World Heritage Site. The park's visitor centre is located at a hill in Namche Bazaar, where a company of the Nepal Royal Army is stationed for protecting the park. The park's southern entrance is a few hundred meters north of Monjo at 2,835 m. Trekking and climbing groups must bring their own fuel to the park (usually butane and kerosene), and the cutting of wood is prohibited. The Sagarmatha Pollution Control, funded by the World Wildlife Fund and the Himalayan Trust, was established in 1991 to help preserve Everest's environment. About a humdred species of birds and more than twenty species of butterflies have made this park their home. Musk deer, wild yak, red panda, snow leopard, Himalayan black bear, Himalayan thars, deer, langur monkeys, hares, mountain foxes, martens, and Himalayan wolves are found in the park
Early expeditions to climb Everest from the Nepalese side started from Jiri. Before the airstrip at Lukla came into existence all the trekking and climbing expeditions to the Everest region started from Jiri. Starting from Jiri, the route passes through the Sherpa villages of the Solu Khumbu, many of them having beautiful Buddhist monasteries.
Lukla, a village in Khumbu boast of the region’s sole airport.Lying at a height of 9000ft, most travelers to this region usually begin and end their adventure in Lukla. The airport was built in 1964 by Sir Edmund Hillary as part of his project in Khumbu region during the early 60s to transport the supplies for the Himalayan Trust projects in the Khumbu region. Today, somewhere between 90-95% of the foreign nationals who reach Lukla, arrive by a half hour flight from Kathmandu.
Namche Bazar is known as the sherpa capital. Namche is actually a village lying at the junction of the Dudh Koshi and a valley that leads to the frontier pass of Nangpa La . It is tucked away in a niche at a height of 7,845 ft. W. H. Tilman and C. Houston were the first westerners to enter it in 1950 and many more have come since then. Facilities like a bank, a post office, hotels and shops where one can purchase climbing equipment as well as tinned food have sprung up over the years. Namche Bazar is the major regional trading center. Its Saturday market or haat is the place where most of the trading takes place. The headquarters of the Sagarmatha National Park is located in Namche.
Thyangboche is famous for the Thyangboche gompa. It is one of the most important centers of Buddhism in the region. The gompa is the largest in the Khumbu region. It was first built in 1923. Destroyed by a fire in 1989, it was rebuilt later on partly with foreign aid. From Thyangboche, one gets a panoramic view of Kwangde, Tawache, Everest, Nuptse, Lhotse, Amadablam, Kangtenga, and Thamserku.
Buddhism is believed to have been introduced in the Khumbu region towards the end of the 17th century by Lama Sange Dorjee. According to the legend, he flew over the Himalayas and landed on a rock at Pangboche and Thyangboche, leaving his footprints embedded on the stone. He is believed to have been responsible for the founding of the first gompas in the Khumbu region, at Pangboche and Thami. Pangboche is the highest year-round settlement in the valley. The Imja Khola, coming from the right, joins the Dudh Koshi River a little above the village. The gompa (monastery) in Phyangboche is thought to be one of the oldest in the Khumbu region.
Khumjung , a village lying west of Thyangboche, is famous for the gompa where the skull of a supposed Yeti, the Abominable Snowman, is preserved under the supervision of the head Lama. The skull seems more like the outer skin of Himalayan Brown Bear, and this is proved by the report of a scientific exploratory expedition conducted by Sir Edmund Hillary, a copy of which is kept in the gompa.
Pheriche is located at an altitude of 13,845 ft. It lies on a level patch. Apart from the basic facilities available here, there is a medical-aid post maintained by the Himalayan Rescue Association of the Tokyo Medical College with Japanese doctors in attendance. Among other facilities, there is an air compression chamber installed for assisting victims of high altitude sickness
The scenic village of Gokyo lies below the hilly Gokyo Ri(5483m). The village is a cluster of stone houses and walled pastures.One has to pass by the holy Gokyo lakes on the way to the village. The Ngozumpa Glacier Nepal’s longest glacier at 25 miles has to be traversed enroute to this remote village. Gokyo Ri looms above the village on the northern edge of the lake. The summits of Everest, Lhotse and Makalu are visible from the summit of Gokyo Ri.
Thami at 3750m is in a large valley. The village has a police checkpost and a few lodges and tea shops. A little above the village is the Thami gompa, which is the site of the annual Mani Rimdu festival.
Sherpas live in the upper regions of Solu Khumbu. They emigrated from Tibet about 600 years ago. In the past they were traders and porters, carrying butter, meat, rice, sugar, and dye from India, and, wool, jewelry, salt Chinese silk and porcelain from Tibet and beyond. The closure of the border between India and China undermined their economy. Fortunately, with the mountaineering expeditions and trekkers, the Sherpa's found their load carrying skills, both on normal treks and high altitudes in great demand. The Khumbu region has provided a strong group of able bodied, hardy and fearless Sherpa porters and guides. The sherpas are Buddhists.
At the lower elevations lives the Kiranti Rai. The villages of Jubing, Kharikhola, Okhaldhunga, are inhabited by the Rais. Of mongoloid stock they speak their own dialect. Reference is made of their fighting spirit in the Hindu epic Mahabharata. The people from this group have supplied recruits to Gurkha regiments both in the British as well as Indian armies. The Rais follow a religion that is partly animistic with a strong Hindu influence. They revere their ancestors by observing Kul or Pitri puja every year.
The Jirels live in the area around Jiri. They are mongoloid and follow Buddhism.
Losar is celebrated in the month of February by the Sherpas. ‘Losar’ means New Year in Tibetan. Apart from the Sherpas and Tibetans, the Gurungs and Tamangs also celebrate Losar. Buddhist monks offer prayers for good health and prosperity at monasteries. People exchange various goods and gifts among them. Families organize feasts and perform dances.
Dumje is celebrated to mark the birthday of Guru Rimpoche (Padmasambhava).The celebration takes place in June and lasts for six days. It is celebrated in a big way in the villages of Namche, Thame and Khumjung.
Mani Rimdu is a festival that celebrates the victory of Buddhism over the ancient animistic religion of Bon. This festival is celebrated in the monasteries of Thyangboche, Chiwang and Thami. At Thangboche the celebration takes place during the November- December full moon. At Thami the Mani rimdu is festival is celebrated during the full moon in May.Chiwang Gompa generally celebrate this festival during autumn. The lamas wear elaborate brocade gowns and papier-mâché masks while performing. Through the dances, symbolic demons are conquered, dispelled, or converted to Dharma Protectors as positive forces clash with those of chaos. The dances convey Buddhist teaching on many levels from the simplest to the most profound, for those who do not have the opportunity to study and meditate extensively. It gives an opportunity to the Sherpas to gather and celebrate together with the monks.
Sakela (Chandi Dance) is a harvest festival celebrated by the Rai community. The harvest ceremony involves the worship of mother earth, called ‘Bhumi-Puja’. The festival is celebrated twice a year, once in spring before planting begins and once during autumn before harvesting. Ubhauli is celebrated during the spring season on Baishakh Purnima. In the autumn season on Mangsir Purnima, Udhauli is celebrated. The spring worship is done to propitiate mother earth for a good harvest and the rain god to bless the earth with enough rain. The festival is celebrated with more fervor in the remote hills. The Rai villagers celebrate it with priests (dhami) who perform rituals to worship their ancestors. The elders of the community begin the dance with a puja. Later on everybody participate in the dance forming a circle by holding each other’s hands. With drumbeats, they begin dancing at a slow pace but moves faster later with the drumbeats. The dance steps and hand gestures imitate the sowing and harvesting of crops .The festival also provides an opportunity for the Rai people to socialise.
The Classic Everest Base Camp Trek
Mt Everest Base Camp is the most popular destination for trekkers in Nepal. Its popularity has grown since the first expedition to the Nepalese side of Everest in the 1950s.One can do this trek the old way, by beginning the trek from Jiri. From Jiri it takes around nine days to reach Namche. On the way you will come across Rai settlements. The other (quicker) alternative is to take a flight to Lukla and to begin the trek from there. The trek follows the Dudh Kosi valley route with an ascent up to the Sherpa capital of Namche Bazaar. From Namche, you traverse along a high path from where you have the first good view of Everest. You head towards Thyangboche Monastery located on top of a mountain ridge and then descend the Imja Khola and continue to the villages of Pangboche and Pheriche. After that you arrive at the Khumbu Glacier. The trek through the glacier takes you first to Lobuche and then to Gorak Shep. From Gorak Shep you can climb up to Kala Pattar for even more spectacular views of the surrounding mountains, including Everest's south west face. Yhou then reach your destination, the Everest Base Camp at the foot of the Khumbu ice fall. |
On Cancer’s Trail
by Florence Williams
Stefanie Raymond-Whish was 9 years old when her grandmother was diagnosed with breast cancer. A traditional Navajo who raised 15 children after her husband died in a car wreck, Raymond-Whish's ama' sa' ni seldom spoke about her illness. Even after her surgery, when she lived with the grandchildren and their mother, she always acted strong around the kids. It became a pattern: When Raymond-Whish was 13, her 38-year-old mother, Nellie Sandoval, was also diagnosed with breast cancer. And Sandoval was equally reserved on the subject. "My mother was really good about not appearing sick in front of us," says Raymond-Whish, now 32. "As a little girl, I knew about cancer, but didn't understand the impact of it at the time."
She understood it better by the time she was in college, in Flagstaff, Ariz., when a new tumor appeared in her mother's other breast. "When my mom had her recurrence, that's when it really hit me ... it was really upsetting. I went home to Farmington for her lumpectomy." Sandoval survived the disease, but not without a long struggle that included chemotherapy, radiation, and finally a double mastectomy. "My breasts were pretty mangled," says Sandoval, now 58. "So I said, 'Just get rid of them.' " Both Sandoval and her daughter have made breast cancer and its impact on Navajos the focus of their lives. Sandoval became an activist and filmmaker, working out of her papaya-colored home in Farmington, N.M. Raymond-Whish has taken her mission a step further: She works as a molecular biologist at the University of Northern Arizona, searching for breast cancer's root causes. "Is there any difference in how breast cancer develops in Native Americans and non-Native Americans?" she asks. One possible - and provocative - answer is emerging from her lab at the university: uranium.
Scientists have long known that uranium damages human cells. But in over six decades of atomic health testing, no one had ever noticed that uranium, at low doses, can act like an estrogen. No one, that is, until recently, when Raymond-Whish and her coworkers observed some unusual effects in lab animals.
Uranium can be found in several of the Jurassic sandstones that lie beneath the Four Corners region like a wrecked layered pastry. The target of frenzied mining throughout the Cold War, uranium ore has been wrenched from the ground, pulverized, milled and tossed in tailings across the Navajo Reservation. Low-level radioactive waste has dissolved into groundwater, escaped onto dust particles and blown off thousands of passing trucks to settle uneasily on surface soils. Over 1,000 abandoned uranium mines pockmark Navajo lands, but only half of them have been reclaimed. Exposure to uranium and its daughter elements has been linked to lung cancer, kidney damage and bone disease in Navajos, and it is the suspected culprit in numerous other medical conditions, from degenerative nerve disease and birth defects to a variety of other cancers.
Raymond-Whish's research lab is tucked inside a neo-Grecian edifice on the Northern Arizona University campus. With her gloved hands in a ventilated booth, the white-coated scientist carefully measures out uranium in solution into small test tubes. The solution will be injected into dishes of cultivated human breast cells, donated by a nun who died of breast cancer in 1979. The MCF-7 cells, as they are known, have been kept alive by the Michigan Cancer Fund through 178 generations of cell division. They are famous among researchers for the properties they exhibit in lab experiments. For example, estrogen causes them to proliferate rapidly - exactly as it does in real-life breast tissue, which is why many women diagnosed with breast cancer have their estrogen-producing ovaries removed. Raymond-Whish wants to see if the cells react in the same way to uranium.
"What I'm really interested in is the development of the mammary gland," says Raymond-Whish, who at this point is just weeks away from finishing her doctoral dissertation. A former teen rodeo star in barrel racing, she once wanted to be a veterinarian. But NAU didn't have a vet school, so she majored in zoology. That eventually landed her in the Discovery Research lab, where she studied the effects of pollution on tadpoles. She found she loved research. "It's like being a detective," she says.
The lab's discoveries have already demolished the conventional wisdom on the properties of uranium. Not only does the heavy metal appear to alter mammary cells at very low doses, but it also seems to interfere with normal hormonal signals. Sometimes the uranium follows the same pathways as estrogen, but sometimes it doesn't, which means it's triggering other endocrine responses as well. "We don't yet know the mechanism of how uranium is affecting these cells," Raymond-Whish says, "but we do know an estrogen receptor is involved. We see it in both animals and MCF-7 cells."
Although the work in Raymond-Whish's lab is considered pure research science, it is impossible to sift it from the real-world context of her family, her culture and her beliefs. Breast, uterine and ovarian cancers have risen steeply in Indian country since the advent of uranium mining. Having watched her grandmother and mother suffer, and now with two kids of her own, 14 and 4 years old, Raymond-Whish can't help but wonder if she's next in line.
But while Raymond-Whish's intimate acquaintance with cancer may harm her credibility as a dispassionate scientist, it may also propel her to help make startling discoveries where no one else has thought to look.
The lab's investigation started several years ago, when Northern Arizona University became part of a team that received a five-year grant from the National Cancer Institute. The project is designed to address community health care, so the local Navajo elders had a few suggestions. They told the scientists they wanted to know more about the health effects of uranium pollution.
"So I started adding uranium to the drinking water of my lab animals," recalls physiologist Cheryl Dyer, who was Raymond-Whish's faculty advisor at the time. "And because I'm an ovarian physiologist, I wanted to see what happened in the ovary." Uranium has long been known to be radioactive and toxic, but no one had ever looked at its effects on follicle counts, or the number of "pre" eggs - eggs in the ovary that have not yet been released for fertilization. Dyer and Raymond-Whish found that the number of pre-eggs declined with low exposure to uranium, and that the mice developed heavier-than-normal uteruses. Normally, a toxic chemical will cause an organ such as a kidney to shrink, not expand. "I said, 'Whoa, what is going on here?' " says Dyer. "I started to wonder if there were other heavy metals that cause these changes, and it turns out cadmium does the same thing. That's when a light bulb went off in my head. Cadmium is an estrogen mimic." All those decades of lab work with atomic elements, and "they had completely missed the boat on estrogen mimicry."
Raymond-Whish was the lead author of a paper showing the unexpected effects of uranium on mouse follicle counts, uterine weights and accelerated puberty. "Drinking Water with Uranium below U.S. EPA Water Standard Causes Estrogen Receptor Dependent Responses in Female Mice" was published in December in Environmental Health Perspectives, a peer-reviewed journal put out by the National Institutes of Environmental Health Sciences. Raymond-Whish concluded that uranium acts as an estrogen, and she recommended that Navajo girls and women be followed closely for reproductive cancers. In conversation, Dyer makes her opinion clear: The U.S. Environmental Protection Agency should lower its drinking water standard for uranium from 30 micrograms per liter to 20 micrograms, the Canadian standard. But Dyer and Raymond-Whish are tip-toeing out on a treacherous scientific limb by suggesting policy changes that are based on controversial data.
Raymond-Whish's work and its results have landed her in the middle of a scientific and regulatory quagmire. It's one thing to regulate a chemical known to be toxic at high doses; it's entirely another to suggest regulating minute levels of a substance that is readily found across a large swath of the American West. Many communities, not just those on the reservation, are affected by uranium. Recent tests in Colorado, for example, revealed that 37 cities and towns in the state depend on drinking water that exceeds federal levels for uranium and its daughter nuclides.
Uranium is not just an emotional issue for Raymond-Whish, but for the tribe as a whole. The legacy of mining the element on the 27,000-square-mile reservation is so deeply and collectively felt that the Navajo Nation banned it altogether in 2005 in the face of globally rising ore prices. During the '40s and the Cold War period, the U.S. government used yellow cake - or milled and concentrated uranium ore - to build nuclear weapons. The government stopped buying the ore for weapons in 1971, but the commercial nuclear energy market picked up the slack until the early 1980s. Only about a quarter of all U.S. uranium miners were Native American - Laguna, Hopi, Zuni and Ute as well as Navajo. But Native Americans have been disproportionately affected: Their tribal lands are still contaminated, and former miners suffer illnesses and deaths for which many families are still awaiting compensation.
Despite the tribal ban, at least five companies are seeking state permits in New Mexico to mine lands just off the reservation, including on tribal allotment land. In Arizona, 700 individual mining claims were filed in 2005. The prehistoric sea and river beds that run underground from Naturita, Colo., to Grants, N.M., and across to Moab, Utah, still hold an estimated 600 million pounds of low-grade ore. But for every 4 pounds of uranium extracted, 996 pounds of slightly radioactive waste is left over, in piles, in pits and eventually in the soil, arroyos and underground aquifers.
Some Western tailings piles, like those outside of Monticello, Utah, or Grand Junction, Colo., have been cleaned up. But those on tribal lands have fallen through yawning bureaucratic and regulatory gaps. It's estimated that up to 25 percent of unregulated water sources on the Navajo Reservation exceed federal drinking water standards for uranium. And many families still haul water from these wells, despite warnings by health providers and advocacy groups.
In her lab experiments, Raymond-Whish applies concentrations of uranium that match those of water supplies in parts of the Four Corners, at or slightly above the current EPA standard. She will treat the mammary cells - which come bathed in a red wash of nutrients that resembles weak Kool-Aid - twice in nine days with differing doses. She will then collect the breast cells, extract their protein signatures, and use a tedious process to examine differences in the number of their estrogen receptors. She will also feed rats different mixtures of uranium-tainted water and examine their mammary glands for altered development. She will compare those results to rats fed a well-known synthetic estrogen, diethylstilbestrol, or DES, and to rats that have drunk plain tap water. She'll look for changes to the mammary glands' terminal end buds, lobules and milk ducts, changes that may make them more prone to breast cancer. The work is controversial, and its implications, both for the science of breast cancer and for the treatment of past and future mining pollution, could be profound.
Like Marie Curie over a century before, Raymond-Whish is both repelled and fascinated by the heavy element's mysterious abilities to alter living cells. In some respects, Raymond-Whish and Curie are not dissimilar. Curie, a Polish Jew working in anti-Semitic France, was the first woman to teach at the Sorbonne. As the first Navajo to be awarded a Ph.D. in the Biology Department of NAU, Raymond-Whish displays a confident ease in navigating a different dominant culture. Like Curie, she is driven by an unrelenting curiosity.
If it was a difficult journey from being Rookie of the Year in barrel racing to creating stunning presentations on heavy metals, Raymond-Whish doesn't show it. She moves through the fluorescent-lit lab in a quiet, deliberate fashion, her long, shiny hair neatly in place.
"I like it that you're working on something no one knows the answers to and you're finding the answers," says Raymond-Whish. She grew up in Colorado and New Mexico with her siblings, stepfather and mother, who was a high school guidance counselor before becoming a breast-cancer activist. Forty-four percent of Navajos do not graduate from high school, but Raymond-Whish's mother made sure that she did. "Everybody's saying it's a big deal for me to get a Ph.D. For me, nothing less was expected than, 'You're going to college.' "
The lab work is routine - even tedious - but it's also demanding and consuming. She is tired. With her oral defense looming before a committee of distinguished faculty, she doesn't slow down. In the mornings, she drops her two kids at school. Her husband, Bryan, a Wichita Indian, works nearby in the university's admissions office. She shuttles from the tissue culture room down a long linoleum-floored corridor to the animal histology lab with its wide-screened computer that magnifies mouse ovary sections 40 times over. Scrolling across the screen to count the follicles is her least favorite job. "I get motion sick," she says. The ovaries dominate the screen like giant pink potato chips, lightly salted.
The science of endocrine disruptors, which studies chemicals that mimic hormones, is a little over 10 years old and still rife with skeptics. It has only been in recent years that very low doses of chemicals - in the parts-per-billion range - have been measurable. (A part per billion is the equivalent of one kernel of corn in a corn-filled silo 45 feet tall.) But natural hormones do their work at these very low levels in the human body. One theory holds that certain environmental chemicals, both natural and man-made, can bind to and deceive the hormone receptors.
These receptors are the signal towers that trigger - or prevent - cellular responses that govern everything from metabolism to sex. Artificial chemicals scramble the signals. They appear to be interfering with normal cellular communication and altering how and when the cells, glands and organs develop. Endocrine disruptors have been implicated in obesity, infertility and the timing of puberty as well as in cancer. When many older women stopped taking synthetic estrogen a few years ago, breast cancer rates in this country dropped for the first time in 40 years. DES, the control substance used by Raymond-Whish, was given to pregnant women to prevent miscarriages up until 1971. Their daughters, who were exposed to it in the womb, have been stricken with unusual reproductive cancers, and recent studies have shown an increased risk of breast cancer as well.
Typical carcinogens cause a cell's DNA to mutate, eventually leading to cancer. Radiation causes the fragile chains of DNA to break, also leading to errors and mutations. Scientists know a lot about these two types of cancer-causing agents. But endocrine disruption is far more mysterious.
Which is why scientists like Raymond-Whish find themselves at a unique moment in science, just as the traditional models of understanding disease are shifting. The field of breast cancer research in particular is driving the debate. Chemicals such as atrazine and DDT (an herbicide and a pesticide, respectively), plastics - such as the bisphenol A compound found in Nalgene that was banned from baby bottles this spring in Canada - and now uranium, are challenging and confounding scientists seeking to understand the actions of chemicals in the human body.
In the dynamic field of environmental health, toxicologists - who study traditional dose-response curves of carcinogens - and endocrinologists - who study extremely low levels of chemicals that do not always follow expected linear curves - frequently disagree. Because it is not yet known exactly how chemicals like uranium act upon cells, some scientists flatly dispute Raymond-Whish's findings. "Uranium is not plausibly linked as an endocrine disruptor," says toxicologist Margaret Ruttenber, director of the environmental health studies program of the Colorado Department of Public Health and the Environment. "There is an absence of a known mechanism."
Louise Canfield, director of the Native American Cancer Research Partnership at the University of Arizona, says: "My personal opinion is that obesity and other lifestyle factors are key risks (for breast cancer), along with access to care. Uranium in drinking water is a health hazard for sure, but I'm not sure it's a primary cause of cancer."
But others consider the work groundbreaking. "This is a science of subtlety," explains Andrea Gore, a neuroendocrinologist at the University of Texas, Austin, and former advisor to the National Science Foundation. "(Dyer's and Raymond-Whish's) work is consistent with other good labs. People criticize the field of endocrine disruption because we don't always understand the mechanisms, but the effects are still real. This is why animal studies are so important. The responses we see in lab animals can happen in humans, because we share the exact same hormones. The estrogen receptor is similar."
Still, more evidence is needed before scientists concede a link between uranium and breast cancer in humans. "You can make a very strong case with animal studies, but it will never be definitive," says cancer expert Joaquin Espinosa, professor of molecular, cellular and developmental biology at the University of Colorado, Boulder. "You hope that nature would have done the experiment for you out there at some point. You need to show that real people are affected."
But epidemiological data on the reservation is hard to come by. For one thing, it's difficult to sort out reliable cancer statistics and their changes over time. Some Navajo elders consult only medicine men, so some cancers go unreported. Cancer itself is translated in Navajo as Lood doo nadziihii, "the sore that does not heal." Some patients do not seek treatment, nor do they even speak of the disease for fear of wishing it upon their families, according to Fran Robinson, a nurse oncologist at San Juan Regional Medical Center in Farmington. Until recently, the Indian Health Services kept haphazard records in which diagnoses went unconfirmed and doctors came and went. For a variety of reasons, including instances of abuse of trust by researchers, the Navajo Nation guards its own data as closely as any member of the former Soviet bloc.
The New Mexico State Tumor Registry keeps statistics on cancers by county, including those on the New Mexico portion of the reservation, which is also where many uranium mines were located. In her published paper, Raymond-Whish cites registry data from the late 1970s showing a 17-fold increase in childhood reproductive cancers there compared to the U.S. as a whole. These are extremely rare cancers that are related to hormone systems. Another study looking at registry data from 1970-1982 showed a 2.5-fold increase in these cancers among all Native Americans in New Mexico. (Although these statistics are not broken down by tribe, most of the Native Americans in the state are Navajo.) A 1981 paper showed a possible link between incidents of birth defects in families and the proximity of those families to uranium mine tailings. The sample sizes of the first two studies were too small to draw solid conclusions, and the birth defects study was flawed, cautions Charles Wiggins, director of the Tumor Registry. He plans to re-examine childhood cancer statistics this fall, using new data gathered since 1982.
Overall, Native Americans in New Mexico actually suffer less cancer than the rest of the country, including about half the rate of breast cancer. But even as breast cancer rates in the U.S. have leveled or dropped slightly in recent years, they continue to grow among Native Americans, and the rate has increased more steeply over the past three decades. Breast cancer is the number-two killer (after heart disease) of Navajo women and the most common cancer found in Navajos. (In the U.S. as a whole, lung cancer is the most common cancer.) Navajos with cancer also suffer higher mortality rates due to poor access to medical care. One study found that between one-third and one-fifth of Navajo breast cancer patients receive substandard care. Relatively more young Navajo women get breast cancer, although much of this can be explained by demographics: Navajos have a younger population than other groups. To the doctors working on the reservation, the anecdotal evidence is disturbing. "When we see women in their 30s with breast cancer, it really knocks everyone for a loop," says physician Tom Drouhard, who has been practicing in Tuba City, Ariz., for 30 years. "Our ladies come in with later stages and higher death rates. It's hard to say what the trends are. All of these tumors are multi-factorial, and uranium could be another thing thrown at it. We are very paranoid about the situation with uranium. We had uncovered tailings five miles from Tuba City for 20 years. It's a reasonable concern."
Two other hormonally active cancers, uterine and ovarian cancers, have doubled or tripled in New Mexico Indians since 1970 while remaining essentially the same for Anglos and Hispanics. But although lung cancer in the Navajo population has been authoritatively linked to uranium exposure, it's harder to make the case for other cancers.
"It's a tough nut to crack," says Wiggins. "The rise in breast cancer everywhere almost certainly has to do with hormones more than anything else. Is something going on with hormones and hormone receptors? Our data is not going to make or break any one hypothesis, because there are a zillion factors going up or down. But you have to take seriously any proposition anyone comes out with, because we just don't have answers yet."
It's difficult to trace a disease to an environmental exposure that may have occurred years earlier. And so far, cancer cases have not been mapped in concert with drinking water sources. "Is there more breast and reproductive cancer here?" asks Dyer. "Yes, but you can't localize it geographically. It would be nice to establish a connection between where people are getting sick and where they drink their water. It's hard to get the data. It's frustrating."
One major effort is currently under way to do just that, but the sickness in question is kidney disease, not cancer. This five-year, $2.5 million study, a collaboration between the University of New Mexico Community Environmental Health Program, the Eastern Navajo Health Board and the Dine Network for Environmental Health, is being funded by U.S. Health and Human Services. The team is compiling illness data from 1,300 Navajos, backed up by urine and blood samples, and then overlaying the results on a map of 160 drinking wells that have been studied for uranium, arsenic and other contaminants. Preliminary data from 550 residents and 100 wells have already shown that living within .8 kilometer of an abandoned mine is a significant predictor of kidney disease and diabetes. Although the science linking uranium with kidney disease is solid, it's never before been demonstrated on a real-life map showing proximity to mines, says Chris Shuey, an environmental scientist at the Southwest Research and Information Center. Once the kidney data are in, the researchers might look at cancer next, he says.
Of course, Navajos are not the only population exposed to uranium. What about breast cancer rates in other areas with better data?
Susan Pinney is an epidemiologist at the University of Cincinnati. She and her colleagues looked at the population surrounding a nuclear processing facility in Fernald, Ohio, which operated between 1952 and 1989. The facility, which made fuel rods for nuclear power plants, was the site of numerous accidental releases of uranium into the surrounding air and water. As a result of a $73 million class action lawsuit in 1990 against National Lead of Ohio and the U.S. Department of Energy, the Fernald Medical Monitoring Project has accumulated 17 years worth of data on illnesses and exposures. Pinney examined the medical records of 8,770 people, including nearly 5,000 women, for a variety of cancers, and was able to model the exposure level of each individual. Her work is still being prepared for publication, and she declined to discuss it. However, a presentation of her preliminary, statistically significant findings last November to the annual conference of the Breast Cancer and the Environment Research Centers is now available on-line. Its provocative conclusion: "For women living within five miles of a uranium processing plant, degree of exposure to uranium particulates was related to risk of incident breast cancer."
A few months ago, Raymond-Whish held a traditional Kinaalda ceremony for her daughter, Darby, to mark her passage into puberty. One of the most important Navajo rituals, it celebrates fertility, the natural order and harmony with the earth through song and prayer. Raymond-Whish's mother was there along with dozens of other relatives, and Dyer from the biology department at NAU also attended. For Raymond-Whish, it was a happy, soulful event, but shadowed by the uncomfortable realities of her career in cancer research. From now until late middle age, Darby will produce the large pulses of estrogen that have been linked to breast and other cancers in so many women. And natural harmony, as her mother knows, is not what it used to be, especially now that pollutants are acting like even more estrogen in our bodies. Raymond-Whish can only hope that Darby's cells have the normal number of receptors, and that her genes and her environment haven't somehow conspired to reprogram her development.
"What does artificial estrogen do to the breast?" she asks. "It depends on the time of exposure. If you look at cells of a younger individual who's not yet through puberty, and you expose them to uranium, then that could promote earlier onset of puberty, earlier breast budding. And if they're exposed in the womb, you could be changing the way the receptors are expressed through life."
In breast cancer research in general, there is a fundamental shift from large epidemiological studies that look at women's current lifestyles and exposures to an examination of what the women were exposed to as children. "Most epidemiology starts with the moment a tumor is diagnosed," said Irma Russo, a molecular biologist at the Fox Chase Cancer Center in Philadelphia. "We need to look at when normal cells may have transformed many years earlier."
Suzanne Fenton, a research biologist at the U.S. Environmental Protection Agency, agrees. "We think one of the main drivers of breast cancer is what changes occurred in very early life to alter breast development. It's a fairly radical re-thinking."
Among other experiments, Raymond-Whish is exposing pregnant rats to uranium in order to track what happens in their offspring. When she tried this earlier with mice, the female pups exposed in the womb entered puberty approximately two days earlier, just as they did when exposed to DES. It's a subtle difference, but when combined with other real-life exposures, it may add up.
Explains Fenton: "It's important to remember that breast cancer risk is likely determined by a number of compounds interweaving with genetic factors and not just any one exposure."
Certainly, lifestyles on the reservation have changed in many ways over the past 50 years. "Once upon a time there was no diabetes here, no diverticulitis, no colon cancer," says physician Drouhard. "We are now exposed to the same things you are: plastics, fast food, obesity. Now everybody I know eats at Kentucky Fried Chicken." One way to learn more is by working with all those nose-twitching rodents in the lab. In the coming months, Raymond-Whish will repeat her experiments, prepare to publish again, and spend more time staring at the nauseating giant ovaries on the computer screen. The rats are euthanized before they actually get sick. Still, she says, they do have to offer up their organs to science. It's not easy for her to kill the animals. "Culturally, it's an issue," she says. "But I'm searching for something that's going to help somebody or even lots of people. I always say, 'Thank you for your life.' "
Florence Williams is a 2007-2008 Scripps Fellow at the
University of Colorado, where she is researching endocrine
disruption and cancer. A former HCN staffer, she currently serves
on the HCN board.
This story was funded by a grant from the McCune Charitable Foundation.
Kathleen Tsosie, who has devoted her life to helping others, now faces the frightening possibility that her breast cancer has returned.
Glenda Rangel and her family grew up drinking from and swimming in water tanks dangerously polluted with uranium.
Nellie Sandoval, the mother of scientist Stefanie Raymond-Whish, has become an outspoken activist as a result of her own struggle with breast cancer.© High Country News |
US 6543011 B1
A method for recording events in Java. According to a preferred embodiment, an automator is attached to a Java applet. Responsive to selection by a user, listeners are added for each event type produced in the Java applet. Each time a specified event occurs, that event is captured and saved to a data structure. The recording of events is performed until the user stops the process.
1. A data processor implemented method of recording events, the steps comprising:
loading an application;
adding listeners for each event type produced in the application;
capturing user generated events;
recording the user generated events to a data structure; and
replaying the user generated events.
2. The method as recited in
3. The method as recited in
4. The method as recited in
5. The method as recited in
6. The method as recited in
7. The method of
8. The method of
9. A computer program product in computer readable media for use in a data processing system for recording events, the computer program product comprising:
first instructions for loading an application;
second instructions for adding listeners for each event type produced in the application;
third instructions for capturing user generated events;
fourth instructions for recording the user generated events to a data structure; and
fifth instructions for replaying the user generated events.
10. The computer program product as recited in
11. The computer program product as recited in
12. The computer program product as recited in
13. The computer program product as recited in
14. The computer program product as recited in
15. The computer program product of
16. The computer program product of
17. A system for recording events, comprising:
means for loading an application;
means for adding listeners for each event type produced in the application; and
means for capturing user generated events;
means for recording the user generated events to a data structure; and
means for replaying the user generated events.
18. The system as recited in
19. The system as recited in
20. The system as recited in
21. The system as recited in
22. The system as recited in
23. The system of
24. The system of
25. A method of recording events, the steps comprising:
loading an application;
adding a listener for an event type;
capturing an event;
recording the event; and
replaying the event.
26. The method as recited in
27. The method as recited in
28. The method as recited in
29. The method of
30. The method of
31. A system for recording events generated by an application, comprising:
a listener; and
a file; wherein
the automator is attached to an application and adds the listener to a component of an application; and
the listener monitors events of the event type received from a system queue within the application and posts received events to the file.
1. Technical Field
The present invention relates generally to computer software and, more specifically, to methods of recording events in Java.
2. Description of Related Art
The evolution of programming languages has, to a great extent, been driven by changes in the hardware being programmed. As hardware has grown faster, cheaper, and more powerful, software has become larger and more complex. The migration from assembly languages to procedural languages, such as C, and to object-oriented languages, such as C++ and Java, was largely driven by a need to manage ever greater complexity—complexity made possible by increasingly powerful hardware.
Today, the progression toward cheaper, faster, and more powerful hardware continues, as does the need for managing increasing software complexity. Building on C and C++, Java helps programmers deal with complexity by rendering impossible certain kinds of bugs that frequently plague C and C++ programmers.
In addition to the increasing capabilities of hardware, there is another fundamental shift taking place that impacts upon software programming, that is the network. As networks interconnect more and more computers and devices, new demands are being placed on software. One of these demands is platform independence.
Java supports platform independence primarily through the creation of the Java Virtual Machine. The Java Virtual Machine is an abstract computer, and its specification defines certain features every Java Virtual Machine must have. However, the specification for the Java Virtual Machine is flexible, enabling it to be implemented on a wide variety of computers, devices, and operating systems. One of the main tasks performed by a Java Virtual Machine is to load class files and execute the bytecodes they contain.
One type of program executed by a Java Virtual Machine is an applet. An applet is a Java program that has a set of standard properties that are defined by the applet class. This class was developed by Sun Microsystems and is included in the standard Java Software Development Kit (Java SDK).
Although, theoretically, a program written in Java for one platform should perform on any Java enabled platform, given the allowable differences among Java platform implementations and other factors, a Java program or applet should be tested on all platforms on which it is anticipated to perform. Since user actions in Java are handled by events, and since it can sometimes take many hours or days for a problem to manifest itself, testing of the entire Java Virtual Machine on a platform can be very tedious. Therefore, it is desirable to provide methods of automating the functional testing of the Java platform on various systems.
However, current methods of automating testing of the Java platform on various systems requires a specialized execution environment, as well as compilation of a separate program. Furthermore, current methods require that the applet or application must be exited before any automation can take place, and they require a significant amount of system resources. Therefore, there is a need for a simpler method of testing the Java platform, that does not require recompilation of code, that does not require the applet or application to be exited before automation, and that uses fewer system resources.
The present invention provides a data processor implemented method for recording events in Java. According to a preferred embodiment, an automator is attached to a Java applet. Responsive to selection by a user, listeners are added for each event type produced in the Java applet. Each time a specified event occurs, that event is captured and saved to a data structure. The recording of events is performed until the user stops the process.
The novel features believed characteristic of the invention are set forth in the appended claims. The invention itself, however, as well as a preferred mode of use, further objectives and advantages thereof, will best be understood by reference to the following detailed description of an illustrative embodiment when read in conjunction with the accompanying drawings, wherein:
FIG. 1 is a pictorial representation of a distributed data processing system;
FIG. 2 is a block diagram of a data processing system that may be implemented as a server;
FIG. 3 is a block diagram of a data processing system;
FIG. 4 is a block diagram of a Java virtual machine (JVM);
FIG. 5 depicts a sample user interface to an applet recorder;
FIG. 6 is a block diagram illustrating how events are currently handled within Java applets;
FIG. 7 is a block diagram illustrating how events are handled when an automator is attached to a Java applet;
FIG. 8 is a flowchart illustrating how the applet recorder functions;
FIG. 9 is a flowchart illustrating the function performed by an automator listener; and
FIG. 10 is a block diagram illustrating the three main parts of an object created by an automator listener.
With reference now to the figures, and in particular with reference to FIG. 1, a pictorial representation of a distributed data processing system is depicted in which the present invention may be implemented.
Distributed data processing system 100 is a network of computers in which the present invention may be implemented. Distributed data processing system 100 contains network 102, which is the medium used to provide communications links between various devices and computers connected within distributed data processing system 100. Network 102 may include permanent connections, such as wire or fiber optic cables, or temporary connections made through telephone connections.
In the depicted example, server 104 is connected to network 102, along with storage unit 106. In addition, clients 108, 110 and 112 are also connected to network 102. These clients, 108, 110 and 112, may be, for example, personal computers or network computers. For purposes of this application, a network computer is any computer coupled to a network, which receives a program or other application from another computer coupled to the network. In the depicted example, server 104 provides data, such as boot files, operating system images and applications, to clients 108-112. Clients 108, 110 and 112 are clients to server 104. Distributed data processing system 100 may include additional servers, clients, and other devices not shown.
In the depicted example, distributed data processing system 100 is the Internet, with network 102 representing a worldwide collection of networks and gateways that use the TCP/IP suite of protocols to communicate with one another. At the heart of the Internet is a backbone of high-speed data communication lines between major nodes or host computers consisting of thousands of commercial, government, education, and other computer systems that route data and messages. Of course, distributed data processing system 100 also may be implemented as a number of different types of networks such as, for example, an intranet or a local area network.
FIG. 1 is intended as an example and not as an architectural limitation for the processes of the present invention.
Referring to FIG. 2, a block diagram of a data processing system which may be implemented as a server, such as server 104 in FIG. 1, is depicted in accordance with the present invention. Data processing system 200 may be a symmetric multiprocessor (SMP) system including a plurality of processors 202 and 204 connected to system bus 206. Alternatively, a single processor system may be employed. Also connected to system bus 206 is memory controller/cache 208, which provides an interface to local memory 209. I/O bus bridge 210 is connected to system bus 206 and provides an interface to I/O bus 212. Memory controller/cache 208 and I/O bus bridge 210 may be integrated as depicted.
Peripheral component interconnect (PCI) bus bridge 214 connected to I/O bus 212 provides an interface to PCI local bus 216. A number of modems 218-220 may be connected to PCI bus 216. Typical PCI bus implementations will support four PCI expansion slots or add-in connectors. Communications links to network computers 108-112 in FIG. 1 may be provided through modem 218 and network adapter 220 connected to PCI local bus 216 through add-in boards.
Additional PCI bus bridges 222 and 224 provide interfaces for additional PCI buses 226 and 228, from which additional modems or network adapters may be supported. In this manner, server 200 allows connections to multiple network computers. A memory mapped graphics adapter 230 and hard disk 232 may also be connected to I/O bus 212 as depicted, either directly or indirectly.
Those of ordinary skill in the art will appreciate that the hardware depicted in FIG. 2 may vary. For example, other peripheral devices, such as optical disk drives and the like, also may be used in addition to or in place of the hardware depicted. The depicted example is not meant to imply architectural limitations with respect to the present invention.
The data processing system depicted in FIG. 2 may be, for example, an IBM RISC/System 6000, a product of International Business Machines Corporation in Armonk, N.Y., running the Advanced Interactive Executive (AIX) operating system.
With reference now to FIG. 3, a block diagram of a data processing system in which the present invention may be implemented is illustrated. Data processing system 300 is an example of a client computer. Data processing system 300 employs a peripheral component interconnect (PCI) local bus architecture. Although the depicted example employs a PCI bus, other bus architectures, such as Micro Channel and ISA, may be used. Processor 302 and main memory 304 are connected to PCI local bus 306 through PCI bridge 308. PCI bridge 308 may also include an integrated memory controller and cache memory for processor 302. Additional connections to PCI local bus 306 may be made through direct component interconnection or through add-in boards. In the depicted example, local area network (LAN) adapter 310, SCSI host bus adapter 312, and expansion bus interface 314 are connected to PCI local bus 306 by direct component connection. In contrast, audio adapter 316, graphics adapter 318, and audio/video adapter (A/V) 319 are connected to PCI local bus 306 by add-in boards inserted into expansion slots. Expansion bus interface 314 provides a connection for a keyboard and mouse adapter 320, modem 322, and additional memory 324. In the depicted example, SCSI host bus adapter 312 provides a connection for hard disk drive 326, tape drive 328, CD-ROM drive 330, and digital versatile disc read only memory drive (DVD-ROM) 332. Typical PCI local bus implementations will support three or four PCI expansion slots or add-in connectors.
An operating system runs on processor 302 and is used to coordinate and provide control of various components within data processing system 300 in FIG. 3. The operating system may be a commercially available operating system, such as OS/2, which is available from International Business Machines Corporation. “OS/2” is a trademark of International Business Machines Corporation. An object oriented programming system, such as Java, may run in conjunction with the operating system, providing calls to the operating system from Java programs or applications executing on data processing system 300. Instructions for the operating system, the object-oriented operating system, and applications or programs are located on a storage device, such as hard disk drive 326, and may be loaded into main memory 304 for execution by processor 302.
Those of ordinary skill in the art will appreciate that the hardware in FIG. 3 may vary depending on the implementation. For example, other peripheral devices, If; such as optical disk drives and the like, may be used in addition to or in place of the hardware depicted in FIG. 3. The depicted example is not meant to imply architectural limitations with respect to the present invention. For example, the processes of the present invention may be applied to multiprocessor data processing systems.
With reference now to FIG. 4, a block diagram of a Java virtual machine (JVM) is depicted in accordance with a preferred embodiment of the present invention. JVM 400 includes a class loader subsystem 402, which is a mechanism for loading types, such as classes and interfaces, given fully qualified names. JVM 400 also contains runtime data areas 404, execution engine 406, native method interface 408, and memory management 424. Execution engine 406 is a mechanism for executing instructions contained in the methods of classes loaded by class loader subsystem 402. Execution engine 406 may be, for example, Java interpreter 412 or just-in-time compiler 410. Native method interface 408 allows access to resources in the underlying operating system. Native method interface 408 may be, for example, a Java native interface.
Runtime data areas 404 contain native method stacks 414, Java stacks 416, PC registers 418, method area 420, and heap 422. These different data areas represent the organization of memory needed by JVM 400 to execute a program.
Java stacks 416 are used to store the state of Java method invocations. When a new thread is launched, the JVM creates a new Java stack for the thread. The JVM performs only two operations directly on Java stacks; it pushes and pops frames. A thread's Java stack stores the state of Java method invocations for the thread. The state of a Java method invocation includes its local variables, the parameters with which it was invoked, its return value, if any, and intermediate calculations.
Java stacks are composed of stack frames. A stack frame contains the state of a single Java method invocation. When a thread invokes a method, the JVM pushes a new frame onto the Java stack of the thread. When the method completes, the JVM pops the frame for that method and discards it. A JVM does not have any registers for holding intermediate values; any Java instruction that requires or produces an intermediate value uses the stack for holding the intermediate values. In this manner, the Java instruction set is well defined for a variety of platform architectures.
PC registers 418 are used to indicate the next instruction to be executed. Each instantiated thread gets its own PC register (program counter) and Java stack. If the thread is executing a JVM method, the value of the PC register indicates the next instruction to execute. If the thread is executing a native method, then the contents of the PC register are undefined.
Native method stacks 414 store the state of invocations of native methods. The state of native method invocations is stored in an implementation-dependent way in native method stacks, registers, or other implementation-dependent memory areas. In some JVM implementations, native method stacks 414 and Java stacks 416 are combined.
Method area 420 contains class data, while heap 422 contains all instantiated objects. The JVM specification strictly defines data types and operations. Most JVM implementations choose to have one method area and one heap, each of which is shared by all threads running inside the JVM. When the JVM loads a class file, it parses information about a type from the binary data contained in the class file. It places this type information into the method area. Each time a class instance or array is created, the memory for the new object is allocated from heap 422. JVM 400 includes an instruction that allocates memory space within the memory for heap 422 but includes no instruction for freeing that space within the memory. In the depicted example, memory management 424 manages memory space within the memory allocated to heap 422. Memory management 424 may include a garbage collector that automatically reclaims memory used by objects that are no longer referenced by an application. Additionally, a garbage collector also may move objects to reduce heap fragmentation.
Turning now to FIG. 5, there is depicted a screen image of user interface 500 for an applet recorder in accordance with the present invention, which may run on top of a JVM such as JVM 400. User interface 500 contains a start record button 510 to start recording events, and a stop record button 520 to stop recording events. User interface 500 also contains a close button 530 to close the applet recorder. The applet is viewed in area 550 on the left side of user interface 500.
The applet is loaded and started prior to receiving a request to record events. Thus, in the embodiment illustrated in FIG. 5, the applet has been loaded and started. Start record button 510 is enabled because recording has not commenced. Stop record button 520 is not enabled, for the same reason.
Turning now to FIG. 6, there is shown a block diagram illustrating normal Java applet operation that runs on top of a JVM such as JVM 400 and may be implemented in a data processing system such as data processing system 300. An applet 620 must be loaded into an applet viewer, such as the applet viewers within Netscape Navigator or Microsoft Internet Explorer. Applet 620 contains all of the user-accessible components. Once applet 620 is loaded, it creates a standard Java class of event listeners (shown in FIG. 6 as applet listeners 640) that are attached to these components and system queue 650. It should be noted that several applet listeners may be (and usually will be) used.
Applet listeners 640 are event listeners. An event listener is any object that implements one or more listener interfaces. There are different listeners for each category of event in Java. For instance, the MouseListener interface defines methods such as MouseClicked, MousePressed, and MouseReleased. In order to receive events from a component, an object adds itself as a listener for that component's events. If an object implements the MouseListener interface, it listens for a component's mouse events by calling addMouseListener on that component. This allows a component's events to be handled without having to create a subclass of the component, and without handling the events in the parent container.
In response to user input 610 on a component in applet 620, such as moving a mouse, a keystroke, or a drag operation, an event 630 is constructed and posted on system queue 650. System queue 650 then dispatches this event to any applet listeners 640 on that component. The component's applet listeners 640 execute tasks according to the properties of event 630. Examples of tasks performed by applet listeners 640 include loading or saving information to a file when a button is depressed, playing a sound or displaying an image when the mouse cursor is moved over a specific area, and closing a program when a specific combination of keys is pressed.
Turning now to FIG. 7, there is a block diagram illustrating an applet recorder 700 in accordance with the present invention. Applet recorder 700 runs on top of a JVM, such as JVM 400, and may be implemented in a data processing system, such as data processing system 600. Applet recorder 700 consists of automator 760, which loads an applet 620 from a database located either on the local data processing system or on a network computer, such as server 104, for viewing. Automator 760 references applet 620 and adds automator listeners 770 to each of the applet 620 components. In response to user input 610 to a component of applet 620, an event 630 is constructed and posted on system queue 650, as is done with normal applet operation as discussed above.
However, system queue 650 not only dispatches event 630 to applet listeners 640, but also it dispatches event 630 to automator listeners 770 on that component. Automator listeners 770 receive an event 630 and store event 630 information to automator queue 780. When the recording session is complete, automator queue 780 contains all of the events that have occurred on applet 620 components. These events can then be played back by being posted to system queue 650 in the same order in which they were recorded.
Automator listeners 770 are similar to applet listeners 640, and are created by the automator and attached to each component of the applet 620. However, rather than perform a specified task to implement applet 620 as applet listeners 640 are programmed to do, automator listeners 770 capture events 630 and record them to automator queue 780, thereby recording the events such that they may be played back at a later time. By having these events stored, testing of a Java Virtual Machine, such as JVM 400, on a particular platform may be automated by having applet recorder 700 replay the user-generated events, thus freeing a person from this tedious task. These events may be required to be played back several times over a period of hours or days.
Turning now to FIG. 8, there is shown a flowchart illustrating a preferred method for recording events with applet recorder 700. After the applet recorder is started, it waits in idle mode (step 805) until it receives an indication from a user to start recording events 630 from an applet 620 (step 810). Once the indication to start recording is received from the user, an applet 620 is loaded into applet recorder 700 (step 815). The applet recorder then adds automator listeners 770 to each component of the applet (step 820), and then waits for user input (step 825).
If user input is received (step 830), then an event is constructed and posted to system queue 650. The system queue dispatches event 630 to automator listeners 770, which capture event 630 (step 835) and record event 630 to automator queue 780 (step 840) for later playback. Applet recorder 700 then continues to wait for user input (step 825).
If no user input is received (step 830), then the applet recorder determines if a stop recording command has been received from the user (step 845). If no stop recording command has been received (step 845), then applet recorder 700 continues to wait for user input (step 825). The recording of events generated by the user through applet 620 thus continues until a stop recording command is received (step 845). When a stop recording command is received, automator listeners 770 are removed, and applet recorder 700 ceases to record events generated by applet 620 (step 850). Applet recorder 700 then idles (step 805), waiting for a command to start recording anew (step 810).
With reference now to FIG. 9, there is shown a flowchart illustrating the function performed by an automator listener 770. Automator listener 770 idles (step 910) until an event occurs (step 920). Once an event, such as MOUSE_CLICKED, occurs, the event information is saved to an object (step 930). The object is then added to the automator queue 780 (step 940) and the automator listener continues to idle (step 910), waiting for the occurrence of another event (step 920).
A block diagram illustrating the three main parts of an object created by an automator listener 770 is depicted in FIG. 10. The object consists of event ID 1010, component field 1020, and event information 1030. Event ID 1010 indicates the type of event that occurred, such as MOUSE_CLICKED, ITEM_STATE_CHANGED, etc. Component field 1020 references the component on which this event occurred, such as a button, list, text area, etc. Finally, each object contains specific event information 1030, which includes several event-specific things.
As an example of the functioning of an automator listener 770 and the creation of an object, suppose applet 620 button is clicked by a user. The automator listener 770 for that component would create an object with MOUSE_CLICKED as event ID 1010. A reference to the button component would be placed in component field 1020. Event information 1030 would contain all other information about the event, such as the x-y coordinate position, modifiers, click count (double or single click), etc.
It is important to note that, while the present invention has been described in terms of recording events generated by Java applets, it is also applicable to applications as well. Furthermore, while described principally with respect to Java, the present invention may be applied to other event-driven object oriented programming languages following a similar “listener interface” model.
It is also important to note that, while the present invention has been described in the context of a fully functioning data processing system, those of ordinary skill in the art will appreciate that the processes of the present invention are capable of being distributed in the form of a computer readable medium of instructions and a variety of forms, and that the present invention applies equally regardless of the particular type of signal bearing media actually used to carry out the distribution. Examples of computer readable media include recordable-type media, such floppy discs, hard disk drives, RAM, and CD-ROMs and transmission-type media, such as digital and analog communications links.
The description of the present invention has been presented for purposes of illustration and description but is not intended to be exhaustive or limited to the invention in the form disclosed. Many modifications and variations will be apparent to those of ordinary skill in the art. The embodiment was chosen and described in order to best explain the principles of the invention and the practical application, and to enable others of ordinary skill in the art to understand the invention for various embodiments with various modifications as are suited to the particular use contemplated.
Citations de brevets
Citations hors brevets |
Conflict and Solidarity in Mudejar Society
The activity of the industrious Mudejar at practically all levels of the economy, and the quotidian mingling of Muslim and Christian in field, workshop, and marketplace provided the foundation of—and the most effective adhesive reinforcing—the structure of Valencia's plural society. A common preoccupation with cultivation and harvest, manufacture and profit fostered among Muslims and Christians a certain homogeneity in outlook and understanding and opened the door to more meaningful social interaction, however hesitant and limited. It is, nevertheless, clear that Valencia's plural social structure was imperfectly articulated and fragile, and that the cement of economic relations was ultimately an unequal match to the solvent of overt religious antagonism. Thus, in 1502 when the Mudejars were faced with the threat of forced conversion in Valencia, many chose flight to the Maghrib with its detrimental material consequences, and a century later, after the abject failure of a policy meant to persuade the Moriscos to become true Christians, the Crown decided that the economic repercussions of the Moriscos' expulsion could be more easily endured than the risks of Morisco dissidence. Both episodes are indicative of the tenacious adherence of Valencia's Muslims to Islam, and this adherence reflects a religious and social reality far more compelling than the material life they shared with Christians. The aim of this chapter will be to provide some understanding of that reality and to explain how the Mudejars maintained both their identity as Muslims and a group cohesiveness while living in a Christian world. It will be suggested that Mudejar self-perception and group-perception were founded not only on religious belief per se, but also on
the perpetuation of a social world the organization and mores of which were distinct from those of Valencian Christians. Indeed, it is perhaps more accurate to view the Mudejars as forming a distinct subsociety separated from the larger Christian society by patterns of intragroup social and cultural behavior, but intersecting with that Christian society in a number of ways not without telling acculturative impact.
The Acculturative Challenge
However durable Mudejar society might have been, it had nevertheless experienced considerable acculturative change. Prior to analyzing late-fifteenth-century Mudejar society, it will be useful to outline briefly some of the more significant of these changes, brought about both subtly by the Muslims' informal interaction with the Christian populace, and forcefully by the formal imposition of Christian authority on an Islamic society.
Valencia's Muslims underwent their most radical cultural adjustment in the half-century immediately following the Christian conquest. During that time the Islamic society was essentially decapitated, losing its political and cultural elites to emigration and in failed rebellions. Islamic government and the public application of the Shariah were rudely shunted aside by a Christian colonial administration with its own body of law, the Furs . Even the autonomous and corporate aljama, vouchsafed to the Muslims by their Christian overlord, had no Islamic antecedent, compelling the Muslim adelantats to administer their communities in a foreign manner analogous to that of the Christian jurates in their corporate universitas . As has been indicated, the Christian bureaucracy, particularly when functioning as a tax-collecting mechanism, intruded on most areas of Mudejar life, so that in even so intimate a family affair as inheritance strict observance of the Shariah had to be modified in compliance with the king's demands. Furthermore, the Muslim by necessity had to learn to function in a Christian as well as Islamic legal system. Although this was to be expected for cases involving both Christian and Muslim parties, the fact that some Muslim litigants pursued civil suits with coreligionists in Christian courts is indicative of an acculturative process potentially damaging to the internal harmony of Mudejar communities.
Also important were the more gradual processes of Christian urbanization and seigneurialization, which, coupled with the decreasing size of the Muslim population (from an overwhelming majority in the thirteenth century to 30 percent of the kingdom's total by the midfifteenth century), altered the very appearance of the kingdom, as churches re-
placed mosques and Muslims retreated from the urban centers to the countryside. Although the rhythms of agricultural and pastoral life remained largely unchanged, Christian control of the economy, especially of the larger domestic and international markets, meant that Muslim farmers, artisans, and retail merchants were dependent on Christians for marketing their produce, selling their manufactured goods, and purchasing raw materials and bulk goods. Mudejar dependence was further manifested in their vassalage to Christian lords and their loss of ultimate control over their land. Outside of the walls of his own home or mosque, the Mudejar could hardly avoid coming to grips with the Christian presence, and having to do so on Christian terms.
The greatest threat to the integrity of the Mudejars' Islamic culture and Muslim identity lay in the cities and towns, where the populations were overwhelmingly Christian. There the expressions of Christian religiosity were most pervasive and most aggressive, while conversely the public observance of Islam was all the more restricted. The best example of this tendency is the anti-Muslim and anti-Jewish preaching of the Inquisitors from the pulpits of Valencia and Cartagena. Another urban danger, particularly in the capital, was the volatility of the Christian mob. Although the evocation of despair fed on fear might have been an efficacious proselytizing tool with some Muslims, it is arguable that the demagoguery of Inquisitors and like-minded clerics and Christian mob violence would have been more likely to repel Muslims from Christianity and to inspire among them a reactionary zeal in their commitment to Islam. The cities posed another threat perhaps more serious than Christian aggression, namely, the possibility of friendly and leisurely contact with Christians and their religion. Adherents of both faiths congregated in city taverns and fonduks, passing their leisure time gambling together, drinking together, and sleeping together. For the Muslim such activities meant a violation of the precepts of his own faith, and signaled his entry into a gray area of cultural amorphousness, where, if his Muslim identity was not overtly challenged, there nevertheless resulted a weakening of his own sense of distinctiveness. That Mudejars flocked to Valencia to enjoy with their Christian friends the spectacle of the Corpus Christi Day processions has a similar significance: it does not mean that these Muslims were on the road to baptism; rather, it indicates that they had acquired non-Islamic cultural accoutrements and forms of behavior born of a long-term exposure to Christian society.
A useful index of acculturation is the extent of Mudejar bilingualism. Burns, in unison with Fuster and Barceló Torres, concludes that in the thirteenth century the Muslim masses were unilingual Arabic-speakers. For midfourteenth-century Valencia Boswell points to only a minimal change in this pattern: "the vast majority of the Mudéjares did not
speak the language of the dominant culture." This was in contrast to the kingdom of Aragon, where the Mudejars spoke Romance but had for the most part lost the ability to function in Arabic. It is thought that throughout the fifteenth and sixteenth centuries this state of affairs remained more or less unchanged. In support of this view historians note the Moriscos' obdurate use of Arabic as their language of daily parlance and the systematic effort of the Christian authorities (from 1565) to discourage the teaching and use of Arabic. Of particular significance is the agreement reached between the Moriscos and Carlos I in 1526, in which the Moriscos maintained that "the greater part of the Moorish men and almost all of the women do not know how to speak aljamia (Romance)." Moreover, the Moriscos stated that they would need forty years in which to learn Romance (Carlos I allowed them only ten). The evidence is impressive, and that both the Mudejars and the Moriscos spoke and read Arabic seems indisputable (see below); however, their avowed unilingualism is open to debate. I would argue for a more extensive Mudejar bilingualism, although with Romance spoken with an imperfect accent and syntax. Given the Moriscos' insincere conversion and their anxiety to forestall thorough catechization in the Catholic faith and Inquisitorial inquiries, their presentation of their linguistic status in the negotiations with Carlos I might have been somewhat disingenuous, a device meant to discourage zealous Romance-speaking clergymen. Clearly, in asking for a forty-year period of grace the Moriscos were playing for time. In the years between 1526 and 1609 there very well might have ensued a decreasing bilingualism. As the Moriscos further withdrew from the cities, their contact with the Christian world would have decreased, while, in a reactionary manner, their desire to avoid Christian culture and to cultivate their own cultural distinctiveness—in effect, to freeze the acculturative process—would have intensified. Just as the Christian authorities realized that the Arabic language had to be removed as an impediment to the effective "Christianization" of the Moriscos, it is equally possible the Moriscos understood that by inculcating their children in the Arabic language, the sacred language of the Qur'an[*] , and by forbidding them any education in Romance they were strengthening their fidelity to Islam. In sum, the Moriscos' linguistic status is not necessarily an accurate indicator of that of their Mudejar predecessors.
The records of cases tried in the court of the bailiff general, in which Muslims appear as litigants and witnesses, are suggestive of a significant Mudejar bilingualism. Normally, if a Muslim witness was in need of an interpreter, a role usually fulfilled by the royal qadi[*] , the latter's participation in the case was explicitly indicated. For instance, Nuzeya, a Muslim prostitute from Oliva, confessed with "Ali Bellvis, qadi ,
intervening." More explicitly, one defendant, Ubaydal Allepus of Bétera, is described as "mal algemiat (i.e., he speaks Romance poorly) ... he does not understand algemia," but his interpreter, the amin[*] Açen Amet, was a "Moor molt algemiat and a person who understands la algemia very well." Yet it is striking that the large majority of the Mudejars appearing in these trial records did not require an interpreter. While it is not surprising that the Muslims of urban morerías spoke Romance, given their constant mixing with the Christian populace, or that the amin of a rural aljama had mastered enough Romance to act as intermediary between the aljama and its lord, it is impressive that many Mudejars from rural areas could testify in Romance. Some examples are Açen Muça of Serra, Homar b. Perellos of Benaguacil, Ubaydal Suleymen of Mirambell, Maymo ben Çabit of Manises, and Abdulcarim of Oliva, among others. Equally significant are the testimonies of Christians about their seemingly routine verbal exchanges with Muslims. Ursula, the daughter of the innkeeper Joan Jeroni, testified how Alasdrach and Abdalla Sinube of Buñol, and Ali Alcayet of Chiva, "were speaking Arabic (alguaravia ) with the said captive Moor," whom they had allegedly helped to escape, and how later, after returning from a meeting with the lord of Carlet, "they requested [from her] a good room for sleeping." Miguel de la Serra, a tailor of Valencia, remembered that on Corpus Christi Day Muslim youths from Chiva had come to his house, eaten there, and then invited him and others to accompany them to the festivities. The odyssey of Angela de Vanya, a prostitute from Cuenca plying her trade in Onda, is revealing. She was first approached by Hadal, a Muslim from Benigazlo (Vall de Uxó), but did not know he was a Muslim, because he spoke "in the Valencian tongue and very suavely and not showing any sign of knowing the Moorish tongue." Hadal compelled Angela to go with him to Tales, a Muslim village (loch de moros ), and brought her to the house of Mahomat Cotalla. There, as the frightened Angela sat weeping in the Cotallas' kitchen, Mahomat and his wife assured her that she would not be harmed. Although these trial records afford only a glimpse at a small cross section of Mudejar life, they do show an erosion of linguistic barriers since the fourteenth century. The increasing size of the Christian population and the variety and frequency of its economic dealings with the Mudejars must have necessitated the latter's acquisition of at least enough Romance to carry on day-to-day affairs. It is unrealistic to assume that Muslim and Christian artisans, farmers, and merchants all had interpreters at their disposal for conducting their mundane but essential business. The documentation does not indicate that this was so. If the Mudejars assiduously cultivated an Arabic culture, there is still no reason to assume that at this point in time, before the mass
conversions, they had any reason for a self-conscious refusal to communicate in Romance. Demography and economy advised otherwise, and linguistic adaptation proceeded apace.
It should not be thought that such acculturation of the Mudejars to Romance-Christian culture as did take place was necessarily a negative or a degenerative process, or a process the flow of which was only unidirectional, although it is likely that the recessive but resilient culture of the Mudejar minority had experienced greater modification from its continual contact with a dynamic and expanding Christian majority. Rather, such acculturation was inevitable, an evolutionary process of adaptation necessary for the social and economic viability of a society that juxtaposed the proponents of mutually hostile ideologies. Even if the behavior of some Mudejars seemed almost to flout Islamic convention, this did not signal a reorientation of their fundamental religious beliefs. (If anything, it attests to a kind of secularization, in which beliefs were not so much altered as ignored.) In a medieval plural society, where identity was finally defined by religious affiliation, cultural erosion and unorthodox conduct were not evidence of Mudejar assimilation into Christian society. Such assimilation—the lessening of social distance, as opposed to acculturation, the lessening of cultural distance —could be achieved only through religious conversion, that is, through a fundamental change of identity. Therefore, in order to grasp the degree to which Mudejar society was truly threatened by absorption into the Christian body, one must broach the question of the extent of Mudejar conversion to Christianity.
Unlike the thirteenth century, with its mendicant preachers, schools of Arabic, and refined techniques of polemic, the result of which had been the conversion of a considerable number of Muslims, the Valencia of the Catholic Monarchs saw no organized ecclesiastical campaign of proselytizing. After recovering from the trauma of the conquest, a factor that probably accounted for much of the Dominicans' early success, the Mudejars had regrouped, and so effectively that the Valencian Church seems to have made little or no further headway in the fourteenth century. This caused the Dominican preacher Vicent Ferrer to lament the inactivity of the clergy in missionary work (1413). The stimulus of Ferrer's zeal led to the establishment of an Arabic school in Valencia during the reign of Alfonso V (ca. 1424), but it seems to have quickly faded into oblivion. As we have seen in chapters 1 and 2, although King Fernando cautiously welcomed the baptism of individual Mudejars, he did not promote conversion on a mass scale. The lords of the kingdom, both ecclesiastical and secular, staunchly opposed any attempts to proselytize their vassals, since they could not wring as much rent out of Christian vassals.
There were compelling social reasons for not converting. The convert, it seems, entered a sort of "no-man's land," being a full member of neither Christian nor Muslim society. If anything, owing to the strong ties of blood uniting Mudejar families (see below), converts tended to associate more with their Muslim relatives than with Christians. The convert Miguel Crestia was implicated along with his Muslim brother Ubaydal Allepus in the murder of another Muslim. A tornadizo of Cocentaina came to the aid of a Muslim relative, a runaway slave from Córdoba. Of course, Mudejars did not regard the conversion of their fellows favorably. Thus, the Inquisition arrested two Muslims who boldly and vociferously attempted to dissuade some Muslims from receiving baptism. Moreover, the convert put in jeopardy whatever inheritances he hoped to receive from Muslim relations. However, it seems that the Crown was no longer confiscating the estates of deceased converts, as had once been the case. Therefore, the potential convert was not burdened by the fear of depriving his descendants of their inheritances.
Even so, Christian society did not welcome the convert with open arms. In previous centuries Christians had insulted converts by calling them "dogs," "renegades," and the like. Given the lukewarm reception of the judeo-conversos by Old Christians and the beginning of an obsession with the purity of Christian blood, it is highly unlikely that the situation would have improved. The murder of Mudejar converts by Old Christians supports this supposition.
Nevertheless, a small number of Muslims chose to convert. The large majority of the proselytes were slaves who had the most to gain from conversion. Slaves with Jewish or Muslim masters had to be manumitted after their baptism, for infidels could not own Christian slaves. One Muslim woman preferred conversion to four years of servitude to a Jew. For slaves with Christian masters, the situation was less clear-cut. The law provided that the slave should be freed only if he converted with the consent of his master. Therefore, the decision to manumit the baptized slave was solely the master's. Some masters granted their Christian slaves freedom, but most did not. Business sense usually prevailed over religious scruples. Runaway slaves sometimes converted so as to conceal their fugitive status. When the Muslim slave Fatima fled from Valencia into Aragon, she became Elinor de Vellasquo. From there she journeyed to Barcelona and to Mallorca, where she was apprehended.
There were some sincere free converts. The best known of these was the son of a faqih[*] of Játiva, who received baptism in 1487, became a priest, and preached to the Muslims of Granada and Aragon. Under the name of "Juan Andres" he translated into Aragonese the Qur'an[*] and the "seven books of the Sunnah ." Of the few other Valencian converts
who appear in the documentation, we know nothing more than their name or place of residence, or sometimes both. Conversion among the Aragonese and Catalan Mudejars seems to have been equally infrequent. In an extremely bizarre incident, a Muslim woman of Albarracín converted just fifteen days after her wedding and then endeavored to retrieve her bridewealth from her Muslim husband. Perhaps conversion was a way of escaping an unhappy arranged marriage, although divorce might have been easier.
Mudejar conversion, then, was not unknown, but it amounted to, at best, a few drops in the font. The large majority had no intention of abandoning Islam.
Mudejar Feuding and Social Structure
A discussion of Mudejar acculturation and religious conversion has served to clarify certain aspects of Muslim-Christian interaction and its impact on the minority culture. It has, however, revealed little about the Mudejar subsociety itself when it did not somehow mesh with the public life of the Christian kingdom. Posing a dichotomy between Valencia's Muslim and Christian societies may appear somewhat artificial, since both groups, after all, lived and labored together in the same kingdom. But occupying the same physical space does not necessitate the rigid conformity of all inhabitants to one all-encompassing social system with its prescribed modes of behavior. Even when dwelling in urban morerías amid large concentrations of Christian population, the Mudejars were enmeshed in distinct networks of social relations articulated in accordance with the moral and material demands of their families.
As to the precise nature of Mudejar social structure, the documentation affords us only precious glimpses. Of course, the Christian bureaucrats had little intrinsic interest in the internal life of the king's Muslim communities. What they have left are details that for them were incidental to their fiscal and administrative concerns, bits and pieces that in the aggregate form a coherent, albeit by no means complete, picture.
An exploration of the royal records for patterns of Mudejar behavior and the social structures underlying them reveals a high incidence of intracommunal violence, greatly exceeding the occurrence of violent conflict between Muslims and Christians (although probably not that between Christians). The discovery of 120 Crown-sponsored truces between antagonistic Mudejar families, not to mention the large number of assaults and murders, suggests that the family feud occupied a central place in Mudejar social life. The Mudejars' internecine violence indi-
cates that their subsociety was still vital, and represents their channeling of energy inward into prescribed forms of behavior and association that had social meaning for them alone. In order to comprehend the origins of such feuding and the historical significance of this phenomenon, it will be necessary first to discuss the Berber and Arab settlement of the Valencian region and its structural implications, and then to describe the structure of the Arabo-Berber family, particularly its Mudejar variant.
Historians differ as to the precise chronology and nature of the Muslim settlement of the Valencian region (Sharq al-Andalus ). Guichard argues in favor of an early (eighth through ninth century) Berberization of the area, but noting that from the eleventh century onward the Berbers were so Arabized that they pretended an affiliation to prestigious Arab tribes, conveniently forgetting their Berber origins. Barceló Torres disputes Guichard's use of toponyms as evidence for early Berber settlement, suggesting that these Berber place names originated in later waves of Berber immigration under al-Hakam[*] II (961–976) and al-Mansur[*] (978–1002), or perhaps during the conquests of the Almoravids and the Almohads. Míkel de Epalza has made important contributions regarding the religious status of Valencia's indigenous population following the Muslim conquest. Taking account of new evidence showing only a limited Christianization of pre-Islamic Valencia, he proposes a rapid mass conversion to Islam (in contrast to Bulliet's model of a more gradual conversion), a trend that was intensified in the tenth and eleventh centuries with a Cordoban-controlled campaign of politicoreligious indoctrination intended to stave off a Fatimid Shic i[*] threat to Sunni[*] al-Andalus. The result was that the Sharq al-Andalus became perhaps the most highly Islamized and Arabized region in the peninsula. From the above it may be said that on the eve of the Christian conquest the Valencian region was Islamized and Arabized with a greater or lesser portion of the population being possessed of Berber roots. Moreover, it is likely that the region was settled on a tribal pattern, in which individual villages were peopled by particular clans. For our purposes, whether these clans were Berber, Arab, or Arabized Berber is not important, inasmuch as Arabs and Berbers had and, indeed, still have, very similar forms of social organization.
The social and political consequences of Arab and Berber settlement, and of Islamization and Arabization have been set forth by Guichard. He convincingly repudiates the views of Sánchez Albornoz and like-minded historians who argue in favor of the assimilation of the Muslim conquerors to the social norms of the indigenous population. According to Guichard, not only did the religion and literate culture of the conquerors become predominant in al-Andalus, but Arabo-Berber tribal
structures prevailed as well. In terms of social organization, this meant that the Arabs and Berbers were members of agnatic, patrilineal groups in which endogamous marriage was preferred so as to maintain the cohesion, wealth, and power of the lineage. These agnatic groups were embedded in a segmentary social system characterized by the balanced opposition between the increasingly inclusive segments of elementary family, lineage, clan, and tribe. This system allowed for both atomization, when segments of equal size competed within the tribe or clan, and amalgamation, when the various segments of a tribe formed a unified front in opposition to an enemy tribe. Thus, tribal in-fighting, most notably between the confederations of Mudari[*] and Yemeni[*] Arabs, marked the political history of al-Andalus well into the tenth century. The intermarriage of Arabs and Berbers with women of the native population—the offspring would have been Muslims and full members of their father's lineage—and the clientage of muwalladun[*] , the descendants of converts to Islam, to particular tribes allowed for the retention of tribal structures.
By the eleventh century tribalism had played itself out as the dominant factor in the political life of al-Andalus. In addition to the political centralization and pacification achieved by c Abd al-Rahman[*] III, Guichard emphasizes the processes of sedentarization and urbanization, which hindered segmentation, the dynamic of tribal organization. David Wasserstein argues that an "Andalusian identity" had supplanted tribalism, any politically significant elements of which were destroyed by the military reforms of al-Mansur[*] , which had dissolved the tribally based jund units. Glick suggests an additional factor—the explosion of conversion to Islam in the midtenth century (according to the Bulliet model) saw the Arabs and Berbers being numerically swamped by neo-Muslims for whom tribal issues were not a major concern. However, if Epalza's thesis of an earlier and more intense Islamization of the Valencian region is correct, it may be that Arabo-Berber structures had taken deeper roots there. This may explain the particular vitality of Arabo-Berber social structures among Valencia's Mudejars.
In any case, the decline of tribalism as a political force, or the dissolution of the tribal unit as a form of social organization, need not have resulted in the lapse of the more elementary segments of agnatic family and lineage as significant structures. In fact, among the Muslims of Nasrid Granada, and even among the Moriscos of postconquest Granada, agnatic solidarity (c asabiyah[*] ) continued to be a potent social force.
It is one thing to remark on the survival of agnatic solidarity in an Islamic polity, or in a recently conquered one like Granada; it is quite another matter to assert that Arabo-Berber social attitudes and structures, at the level of family and lineage, remained largely intact in
Valencia some 250 years after the submission of that region to western Christian domination. The defeat and emigration of Muslim military elites in the thirteenth century had eliminated whatever political significance still remained to either the tribe or clan. Since the government was Christian and denied Muslims access to political power, it is difficult to see what sociopolitical role and aims even the smaller solidarities of family and lineage would have had. Furthermore, progressive Mudejar emigration to Granada or the Maghrib, coupled with the internal migration and fragmentation brought about by seigneurial alteration of rural settlement patterns, must have severely strained or severed the ties binding some Muslim lineages together, in effect, fostering a more radical segmentation.
Still, at the same time, the conquerors' granting of communal and legal autonomy to the Mudejars was probably conducive to the retention of preconquest structures, or at least to the alleviation of the pressures for fragmentation. As will be seen below, the various official posts of the corporate aljama, however much they were distortions of their Islamic antecedents, became sources of prestige and power within the Mudejar community, prizes over which rival families deemed it worthwhile to feud. Also, Maliki[*] law, according to which the Mudejars regulated their daily lives, seemed to justify the perpetuation of feuding relations by its provision for either retaliation or the payment of blood money (diyah ) in cases of homicide and assault.
Aside from the persisting legal and institutional supports for traditional patterns of social relations, it is important to emphasize the peculiar resilience and structural coherence of medieval Islamic society at its lower levels, despite the frequent and sometimes traumatic political upheavals in its upper reaches. In the Valencian case, this resilience derived from the ability of the larger solidarities of an Arabo-Berber society to segment without disturbing the primary social structures. Thus, if a family were broken up—say, through the emigration of one of two brothers to the Maghrib—the sons of the remaining brother could create a new agnatic solidarity composed of themselves, their father, and their own male children, even though their position would be weakened through the loss of their paternal uncle's support. In other words, so long as traditional attitudes were not eroded, preconquest structures could be perpetuated merely through biological reproduction. If 250 years seems an incredibly long time for such structures to have endured, it was also more than enough time for them to have reemerged and solidified. Since Mudejar political life was reduced to the confines of the local morería or village there was no impetus for the amalgamation of clans or tribes; family and lineage, however, remained important.
That lineage long continued to be a special concern of Valencian Muslims is evinced in the claims of a Morisco family in 1567 to have descended from al-Mansur[*] and Valencia's Muslim rulers.c Asabiyah[*] in Valencia was buttressed through the practice of endogamy, that is, through parallel cousin marriage with the daughter of the paternal uncle (bint al-c amm). The negotiations of the newly converted Moriscos with Carlos I in 1526 leave little doubt that this type of marriage had been the preferred one among the Mudejars:
inasmuch as among the Moors today there are many marriages concluded between close relations in a degree prohibited by the Christian law and permitted by the Moorish law, which law permits marriage to the degree of that between cousins-germane—the children of two brothers inclusive—should the said marriages begin to be disturbed, and to prohibit those marriages which could be made from today henceforth, would result in the greatest damage and disturbance among the said Moors.
The Moriscos then requested that Carlos intercede with the papal legate and persuade him to grant a dispensation for the endogamous marriages consummated before the mass conversions.
Our documentation is less useful on this matter. Normally, only the name (ism ) of the wife is given—as Fatima, wife of so and so—without ascription to her lineage (nisbah ). This was in keeping with the far greater importance given to patrilineal descent, whereby the offspring of the marriage pertained to the lineage of the father. Of course, in the instance of marriage with the bint al- c amm, the husband and wife would have the same nisbah . In a few rare cases the documentation shows unequivocally the marriage of parallel cousins. In Alcira, Fotoix bint (daughter of) Mahomat Xativi was the betrothed of Ali ibn (son of) Abrahim Xativi. Two cases from Aragon (where the practice of endogamy is all the more impressive, given the greater acculturation of Aragonese Mudejars) suggest that although the right of the male cousin to the hand of the bint al-c amm was recognized, the prospective brides were sometimes less than enthusiastic about their fate. For five years there was pending in the courts of the kingdom a decision on the litigation between a reluctant bride, Fatima bint Mahomat Margnan of Huesca, and her cousin, Ybraym ibn Mofferiz Margnan, concerning the marriage that Ybraym and his father "were claiming" to contract with her. Fatima bint Jayel de Gali of Zaragoza asked for a divorce from her cousin, Faraig ibn Juçe de Gali, on the grounds that he had maltreated her. Faraig maintained that his mother-in-law had put Fatima up to it and pleaded before the king that his wife be restored to him.
Questions of marital harmony aside, there were compelling reasons for the practice of endogamy. The marriage of paternal cousins en-
hanced the solidarity of the lineage by linking the interests of brothers and their children, and giving the bride's kin more control over the bridegroom. They could then present a more powerful front to rival lineages. Moreover, the reproductive power of the daughter was retained within the lineage; more children, particularly more sons, increased prestige, power, and economic potential. The economic angle was crucial. Since among the Mudejars daughters were able to inherit from their fathers, the exogamous marriage of a daughter implied the loss of property by her lineage. True, the groom had to render a specified bridewealth (sadaq[ *] , Catalan accidach ) to his intended, but, because that was paid to and retained by the bride, her family did not profit from it. In some instances wives left their husbands and took their bridewealth home with them, or when couples separated the wives demanded their bridewealth, but it is highly unlikely that families planned for such eventualities.
While endogamy was preferred, exogamous marriage was not at all unusual. Demography militated against every person being able to wed his or her cousin. More important, exogamy served to define relationships between various lineages. Marriage ties created a community of interest, mitigated conflict between families, and stabilized community life. That Mudejars regarded such affinal connections as determinants of conduct is demonstrated in the judicial appeal of Çaat Siquuti. The lieutenant bailiff of Játiva, with the counsel of the local qadi[*], Yuçeff Alçamba, had passed sentence against Çaat, which sentence was then upheld by the bailiff general, with the counsel of Mahomat Bellvis, the qadi general. Çaat complained that the final ruling had gone against him because the two qadi s "are joined by a certain affinity, that the son of the said qadi Bellvis had married the niece of the aforesaid Yuçeff Alçamba, daughter of Yahye Alçamba his brother." In another case, Mahomat Perpir, when he fled to the Maghrib in 1501, was recorded as having Abdalla Murçi and Yuçeff Zignell as brothers-in-law. In 1499 he had reached a truce with Azmet Murçi, a relation of Abdalla, and in 1500 had done the same with Yuçeff Zignell. It may be that affinal ties had helped to close the rift separating the Perpirs from the Murçis and the Zignells, or perhaps in the end they proved insufficient and Mahomat had reasons other than the fear of forced baptism for taking refuge in the Maghrib.
Since endogamous marriage was practiced, and the bonds between agnates thereby strengthened, it should follow that agnates acted jointly in the conduct of feuding relations. This is precisely the state of affairs encountered in the documentation. In a number of cases the perpetrators of violence were two or more agnates. Homaymat and Çuleymen Montani of Alacuás, with accomplices, shot and killed with a crossbow
another Muslim vassal of the seigneury. Ali Orfayçi and his brother Mahomat Çaffahi of Alcira burst into the house of Pedro Delgado and wounded Juçeff Bolaix with two lacerations. On different occasions Muslims of Valldigna were wounded by Mahomat and Abdalla Giber, and by Çat and Ali Bolarif. Mahomat Malich had the misfortune somehow to incur the wrath of Ali and Mahomat Guayna of Artesa. Clearly casabiyah[*] was crucial for providing the strength in numbers that allowed for such aggression.
At the same time, the readiness to resort to violence, combined with the support of sons, brothers, uncles, and cousins, tended to discourage the violent initiatives or retaliation of rivals. The threat posed by the enemy's group solidarity determined the very nature of much of the violence that occurred. Assault and murder were often committed under the cover of night; ambushes were laid and enemies attacked at unexpected moments; and victims were often alone and outnumbered by their assailants who struck in tandem with agnates or other—perhaps affinal—accomplices. We have mentioned elsewhere Juçeff Cabot of Játiva who returned to Valldigna, where, one evening, with henchmen, he did away with an enemy in his own home. Abdalla Çentido and Fuçey Ylel dispatched Azmet Gradano with a dagger-thrust to the throat as he was leaving his father's house in Mirambell. Mahomat Flori of Játiva was fined 110s, in lieu of eighty lashes, "because he had hidden in his house certain Moors [relatives] from Gandía and from the Vall d'Alfandech [Valldigna] for the purpose of killing Gabix, Moor of the said morería ." Mahomat Chiquet of Alcira even resorted to arson, attempting to burn down the house of the Arrayço family. Only rarely were attempts made on the lives of more than one member of a family at the same time, and in such cases the elements of surprise, darkness of night, and superiority in arms are apparent. The brothers Çahat and Amet Pachando were the victims of a ferocious nocturnal assault in Liria by Caet Natjar and Abraym Rabaça of Bétera. Natjar practically cut off Çahat's head with his sword, while Rabaça seriously wounded Amet with a poisoned crossbow dart.
As casabiyah[*] was predicated on the responsibility of each agnate to uphold the honor of his family, an attack on one member of the group demanded of the others reprisal against the offender or his agnates. Failure to fulfill this responsibility resulted in the family's loss of honor and in a decrease in its prestige and power in the community. Thus violence elicited a violent response, setting in motion a potentially endless sequence of aggressive acts characteristic of the feud.
The understanding of family members involved in feuds that violence would inevitably ensue and vengeance be exacted prompted Muslims of Oliva, Carlet, Játiva, the Vall de Uxó, and Valencia to apply for licenses
to bear arms for the purpose of self-defense against their "enemies." Yuçeff Albanne's application hinged on the presumption that because Abrahim Corumbell had already wounded his brother he was next in line.
A consideration of the feud between the Murçi and the Torralbi families of the morería of Valencia demonstrates the centrality of agnatic solidarity in the prosecution of a feud. We first encounter the brothers Azmet and Çahat Torralbi receiving from the bailiff general arms-bearing licenses (20 March 1503). Because Çahat had denounced Azmet Murçi before the bailiff's court (for a reason unknown to us), he believed he needed to be armed so that he would not be "damaged by the said Azmet Murçi or by another relative of his " (italics mine). Azmet Torralbi, sharing responsibility for his brother's legal actions, whether he liked it or not, felt equally threatened. Such was the expectation of Murçi vengeance that the Torralbis were licensed to bear arms not only for self-defense but also for the defense of "one other companion who will go in your company." Clearly, it was inadvisable for a man embroiled in a Mudejar feud to walk alone.
Less than three months later the Torralbis took the offensive and, putting his licensed arms to good use, Azmet Torralbi severely wounded Abrahim Murçi. The precise relation of Abrahim to Azmet Murçi is unknown, although it seems most likely that they were brothers. In any case, they had the same enemies. Both Abrahim and Azmet had concluded truces with Mahomat Perpir and Ali Perpir, and both had fallen out with even the same agnate, Ubaydal Murçi, perhaps their cousin. Abrahim had also made Azmet's dispute with the Torralbis his own. Nusa, Abrahim's wife, maintained that Azmet and Çahat Torralbi, and their father Abdalla, all bore a grudge against her husband for no reason, and that Çahat and Abdalla prompted Azmet and planned with him the attack on Abrahim. She also asserted that as early as the previous January (1503) Çahat Torralbi had shown much bravado toward her husband and an intention to kill him.
The testimonies of other witnesses contradicted some of Nusa's charges. First, it seems that Abrahim Murçi was not entirely blameless. Abrahim was said to have a "foul mouth," and many thought that "because of his tongue the said Azmet Torralbi has wounded him." Abrahim had indeed gone about stating that Azmet was a fool and deserved to be treated as such. Second, it seems that Çahat and Abdalla Torralbi had not urged on Azmet or planned the crime with him; rather, the attack was more of a spontaneous act on Azmet's part. In fact, Abdalla Torralbi had attempted to mollify his son's animosity toward Abrahim, urging Azmet to make peace with him. On the very day of the assault, Abdalla confronted Abrahim, demanding to know why he was
deliberately antagonizing his son. Abrahim retorted that he did not have to answer to Abdalla, then wheeled and strode off.
Abrahim's inability to hold his tongue proved calamitous, for when Azmet vented his wrath against him, he showed little mercy. Azmet slashed Abrahim about the head and arms with his sword, and, when Abrahim fled into the home of Gil Sanchiz, he followed him there and cut off his right foot. According to the painter Gabriel Gosalbo, it was his own intervention that prevented the affair from escalating into a bloodbath. With a lance he barred the way into Abdalla Torralbi's house against those (Abrahim's relatives, perhaps) who spoke of entering to slit the throats of Abdalla's daughters. Later, he met Çahat Torralbi and advised him to avoid the area of his father's house, lest he fall victim to Murçi revenge.
In another case, the brothers Ali and Azmet Thorruc of Millena killed Azmet Araye of Benilloba on account of some disagreement. Consequently, Azmet's brothers, Mahomat and Çahat Araye—their father was an eighty-year-old invalid—recruited some neighbors in Benilloba and proceeded to Millena, where they took their revenge, slaying Azmet Thorruc.
Açen Muça of Serra was similarly motivated by a desire to vindicate the murder of his half-brother Azmet Gradano ("jermans de mare," a blood tie not as strong as that between sons of the same father) by Abdalla Çentido. For four years he had restrained himself, but on Corpus Christi Day, 1491, he stabbed Abdalla to death as he sat watching the processions in Valencia. Earlier that day, Yuçeff Ada had asked Abdalla about the state of his relations with Açen Muça, and, upon learning that the two were still feuding, admonished Abdalla to be careful. The feud, then, seems to have been a matter of (Mudejar) public knowledge. This makes perfect sense, for the honor and status of a family rested on public evaluation and approval. This approval could best be attained through the display of solidarity and the willingness of the agnates to fulfill their responsibilities to the group, that is, through the family conducting itself properly in a feud.
A number of victimized Mudejar families, either too weak in numbers or too law-abiding to retaliate against their enemies in kind, or perhaps exceedingly confident in the efficacy of royal and seigneurial justice, beseeched their Christian overlords to mete out the punishment due for assault or homicide. The parents of the murdered Azmet Zichnel of Valldigna, "although they have the ability and it would be permitted to them to kill with impunity Juçef Cuyta, the killer," turned first to the criminal justice of Valldigna, a seigneury, and then, once Juçef had fled the valley, to the Crown "to have him fittingly punished and chastized through the measures and means of that [justice]."
It is due to this type of legal recourse to the Christian authorities that much of the information concerning Mudejar violence turns up in the documentation.
The records of the kingdom's fiscal auditor (Maestre Racional ) indicate that royal officials, and probably their seigneurial counterparts as well, acted as mediators between feuding Mudejar families. The result of a bailiff's intervention in a dispute appears in the records as the payment of a monetary settlement (composicio ) by the offending Muslim or his family to the bailiff. Of course, payment to the Christian official provided little compensation for the victim and his family. It is probably safe to assert, therefore, that either the bailiff and the victim divided the settlement or the victim received a separate settlement equal to that received by the bailiff, who, in either case, would have presided over the entire transaction. This conjecture is substantiated in some cases by scribal notations that state that the victim "admitted" or agreed to the settlement. In the eyes of the victim and his family, the settlement would have appeared a form of blood money releasing them from the obligation, although not necessarily from the desire, to retaliate violently.
A more effective means of maintaining the public peace was in the hands of the bailiff general, namely, the official "peace and truce," supposed to last for 101 years, or its permanent variant, the "final peace." Although many of these officially sponsored truces involved Muslims of the morería of Valencia, for keeping order was most difficult in the capital, many others bound Muslims from all over the kingdom. The truces highlight the importance of agnatic solidarity by emphasizing that each party swearing to abide by the truce was representing himself and his "relatives, friends, and defenders." In some cases entire families were present at the conclusion of the truce. Unfortunately, only in rare instances do the documents state the causes for the mutual hostility necessitating a formal truce. When given, the causes usually are acts of violence. The texts of the truces are written according to standard format, in which steep monetary fines are emphasized as the deterrent to any resumption of the feud.
For example, in a truce dated 31 March 1489 we find various members of the Capo family of Alcira, standing "for themselves, friends, relations, and defenders," and Abdalla Pollet, his son Mahomat, Azmet Biari, Çahat Biari and his son Caleth, and Çapdon Eça, all of the same morería, pledging "a good peace and truce between them, to last for 101 years, concerning whatever debates, quarrels, rumors, ill-will, and wounds that might have existed between them until the present day." By virtue of an oath sworn "with hands and with mouth" in the presence of Joan Aduart, royal constable and vicar of the bailiff general, "'to our
Lord God and to the qiblah of Muhammad' the face turned toward midday, according to the Sunnah and the Shariah of the Sarracens," both parties promised that "there will not be done nor caused to be done nor arranged nor attempted, neither openly nor secretly, neither directly nor indirectly, any evil or damage on the persons or goods of them [the other party]." For any violation of the truce, the Crown would penalize the offending family with a fine of 500 florins—100 florins to the injured party and 400 florins to the royal treasury.
It is not difficult to understand why the royal authorities expended such effort to curb the Mudejars' feuding propensities. The death of vassals, the destruction of property, and the disruption of economic activity, all of which resulted from feuding, had the most negative implications for the state of the royal treasury. Moreover, the control of social violence and the efficient administration of justice were essential components of the effective exercise of royal authority. Stated more simply, public order had to be maintained. Feuds not only affected individual families but also at times threatened to engulf entire morerías in tumultuous and bloody disorder. Three royal aljamas were seen to be tottering on the edge of such calamity. In Valencia, the feud between the Roget and the Bizquey families had become so serious that they and "their fathers and mothers" had to be banished from the morería . In Alcira, the bailiff believed the aforementioned truce between the Capo and Biari families to be necessary in order to "pacify the said morería and to put the [morería ] in repose." The bailiff general urged the bailiff of Castellón de la Plana to do something to resolve the conflict between the Bocayos and their rivals "for the benefit and repose of that morería, which today ... is on the road to destruction."
While official intervention prevented the growth of widespread violence within morerías, it is less certain that the truces between families always had their intended durability. Social anthropologies have observed in Mediterranean and Middle Eastern societies where the feud is a central element in social relations the ultimate inefficacy of truces and various forms of compensation as means of permanently extinguishing a feud. Rather, despite the fact that hostilities cease for varying lengths of time, feuds tend to be perpetual in nature. The Valencian evidence does not contradict these conclusions. Azmet Coxet of Paterna was to be apprehended "for breaking the peace and truce." Although Azmet Aixbir and Çaat Borrabe had concluded a truce, Azmet later seized the opportunity to attack Çaat. One also encounters the repetition of truces between the same families, such as the Murçis and the Perpirs, or the Bizqueys and the Rogets. Even if different members or combinations of members of the families appear in successive truces,
this hardly masks the fact of the continuing and potentially explosive state of animosity existing between them.
It should not be thought, however, that without the intervention of royal and seigneurial officials Mudejar society would have destroyed itself through unabated internecine violence. The Mudejars had their own mechanisms for achieving a cessation of hostilities, a state of affairs necessary if routine social and economic life were to continue.
A settlement arrived at through the efforts of Muslim jurists preceded some of the truces established through the office of the bailiff general. When the powerful Paziar family of Alcira and the Getdi family of Picasent concluded an official "peace and truce" after the killing of Abdalla Getdi, the Paziars came with a "carta morischa (Arabic letter) received from Ali Bellvis, son of Mahomat Bellvis, qadi[ *] of the lord king," while the Getdis had an Arabic letter from their local faqih[*] . In other cases Muslims abandoned legal initiatives against enemies, having come to an understanding with them by their own methods and for their own reasons.
In the feud between the Araye and Thorruc families, rather than pursuing the prosecution of Mahomat and Çahat Araye in the governor's court, the Thorrucs dropped the charges against them, after having "established peace with all [their] adversaries." It may be inferred that the Thorrucs realized that the capital punishment of the Araye brothers, even if executed through proper legal procedure, would have served only to provoke the perpetration of retaliatory violence by the Araye and their friends in Benilloba. Each family had already lost a son; with the score even, an uneasy truce seemed wiser than more killing and the disruption of the activities of those still alive. When the wife of the murdered Ali Dabbau, his sister, and the guardian of his children, Azmet Pulpul (apparently not an agnate of Ali), dropped charges against the killer, Çaat Melich, the reasons for their doing so were even more pressing than those motivating the Thorrucs. Here, the surviving members of the victim's family, two women and young children, were incapable of either exacting revenge themselves or defending themselves against later retaliation should they press prosecution through the bailiff's court. This evidence suggests that what moved Mudejars to accept compensation or blood money instead of physical retribution, or to make peace with the enemy, however temporary that peace might be, was not so much the threat of censure by the Christian authorities as the fear of triggering further violence, of reactivating the feud in its most destructive form. In other words, the most effective deterrent to feuding was the feud itself. Rival families eyed each other warily and exercised restraint, committing violent acts sporadically in an
often calculated manner. It was to no one's interest to give violence free rein.
The plethora of official truces between rival Mudejar families, the numerous acts of violence in which agnates were implicated, and the importance of vengeance as a motive for such violence all indicate that the feud was so pervasive as to constitute a primary determinant of Mudejar social relations. Jacob Black-Michaud, whose conclusions are based on the studies of various Mediterranean and Middle Eastern feuding societies, goes so far as to state that "feud can be regarded as a social system per se." While it is not our intention to discuss the validity of this conclusion, nor, for that matter, to attempt to fit Mudejar feuding within the framework of an anthropological model, it is crucial to emphasize a viewpoint on which most observers of feuding societies would seem to agree, namely, that feuding is better viewed as a social process than as a social aberration. If this is the case, then there must have been something other than the necessity to respond to a previous act of violence behind much of the Mudejar violence we have observed, a stake, or stakes, worth the risk of initiating a feud. These stakes were wealth and power. Among Mudejar farmers and artisans wealth was attained through the acquisition of land and through the control of a limited market for manufactured goods. Power, or local influence, rested on the prestige and status afforded by the possession of wealth and by the defeat of rivals in the competition for it. The agnates who constituted the feuding group also functioned as an economic unit, holding land jointly and practicing the same crafts, so that mutual material concerns strengthened agnatic solidarity. Mudejar feuding, then, may be interpreted as a consequence of the competition for material wealth and local status, and as a process determining the allocation of these scarce commodities, thereby stratifying individual communities.
The official truces, which sometimes indicate the professions of the subscribing parties, strongly suggest that economic competition was at the root of much Mudejar discord. In a number of cases both of the feuding parties were practitioners of the same craft, producing for the same limited market. Conflict occurred between Muslim blacksmiths, shieldmakers, and shoemakers in Valencia, between Muslim fishermen in Oliva, and between hemp sandalmakers in the Vall de Uxó. The competition seems to have been most intense among Valencia's Mudejar shoemakers. The feuding Bizquey and Roget families, whom the authorities had to expel from the morería, both practiced shoemaking. The Bizqueys also concluded truces with other enemy shoemakers, namely, Çatdon Caeli, the brothers Abrahim and Çalema Cabero, Çilim Maymo, and Çahat Perpir with his nephews Ali and Mahomat Perpir. The Perpirs themselves clashed with shoemakers other than
the Bizqueys: Çahat Carcaix, Azmet Murçi, Ali Maguarell, and Çaat Abducarim. In the feud which resulted in the wounding of Abrahim Murci, a shieldmaker, by Azmet Torralbi, a shoemaker, it is interesting that the Torralbis seem to have first come into conflict, not with Abrahim, but with Azmet Murçi, a shoemaker (see above).
However, in many other feuds the parties were not of the same profession, in which cases one cannot delineate so precisely the clash of competing economic interests. The ties and common interests created through exogamous marriage or simple friendship complicated intracommunal relations considerably, so that families came into conflict who, had they been guided by economic concerns alone, otherwise might not have. Thus, one encounters truces like the one concluded between the shoemakers Ali Perpir of Valencia and Azmet Naixe of Mislata, on one side, and Ali Maguarell and Çaat Abducarim, shoemakers, and Azmet Claret, a linen salesman, all of Valencia, on the other side. The documentation does not reveal what, other than shoemaking, brought together Perpir and Naixe, Muslims from different families and locales, or why a linen salesman was involved in the feud at all.
The frequent and varied business transactions between Mudejars of all walks of life, and matters associated with the complex pattern of land tenure, in which Muslim artisans also were concerned, provided ample opportunity for the sparking of controversy and mutual hostility. For instance, Çale Magarell had Çaat Feçi imprisoned for money Çaat owed him for the purchase of a donkey. In Játiva, Yuçeff Redona complained that Azmet Beniale had planted mulberry trees on his land and demanded that they be uprooted. Muslims of Valldigna went to court over the alleged sale of cloth, while Muslims of Alcira disputed the ownership of goats.
Material interests fueled the fires of dissension even within families. For instance, Ali Gehini, a wealthy amin[*] of la Foyeta, was so afraid of being robbed by his own sons, disreputable characters who frequented taverns, that he hid his money in the walls of his house. Mahomat Negral had the justice of Valldigna sell his brother Abducalem's mule for debts Abducalem owed him on account of the justice's earlier sale of Mahomat's field for debts that had really been Abducalem's. Unfortunately, because most Mudejar civil litigation was handled in Islamic courts, the records of which do not survive, our information on Mudejar disputes over land or other commodities is extremely limited. Consequently, direct correlations cannot be made between such disputes and the occurrence of violence and feuding. Considering the extent of feuding between Muslims of the same profession reflected in the truces and the few instances of Mudejar property litigation encountered, it can be
cautiously postulated that much of Mudejar feuding had its roots in conflicting economic interests.
The feuding that threatened to destroy the morería of Castellón de la Plana was the violent manifestation of a struggle for political power. The struggle centered on the control of the two posts of adelantat and pitted Abdulazis and his son Yuçeff Bocayo against a faction headed by Çale Arroçen and Ubequer Faraig, which also included Caet Fando, Yuçeff Salio, Ali Gordo, and Ali Gerret. Shortly before 19 April 1487 the Bocayos informed the bailiff general of a fight that had broken out between Yuçeff Bocayo and Ali Gerret. Within weeks matters took a more serious turn: two Christians and some unidentified Muslims entered the Bocayos' home and wounded Yuçeff in the arm and hand, cutting off his finger. Investigation revealed that this was not random violence, but that the two Christian assailants had been hired by Çale Arroçen and Ubequer Faraig to do the dirty work. At this point the bailiff general had Ali Bellvis, the qadi[*] general, intervene. Bellvis managed to persuade the Bocayos and their rivals to agree to a power-sharing arrangement. According to this arrangement, the ten councillors of the aljama, among whom were members of both feuding factions, and the adelantats from the previous year would elect the amin[*] and the two adelantats . Most important, it was stipulated that "in any year either Ubequer Faraig or Çale Arroçen is elected as jurate [adelantat ] that in such case let there be elected as jurate one of the said Bocayos, either the father or the son." That the two posts of adelantat were so hotly contested was probably due to the adelantats ' function as advisors to the amin in the apportionment and collection of taxes. This is suggested by the fact that Ali Bellvis also inspected the aljama's account books to ensure that the amin had justly confiscated certain goods of the Bocayos and Yuçeff Pollina, presumably for reason of unpaid taxes. The Bocayos, then, had a clear material interest in being elected as one of the adelantats . Eventually (August 1488), the authorities established a formal truce between the two factions; however, by 1492 Yuçeff Bocayo was again complaining about wounds inflicted by his former antagonists (it is not clear that these were new wounds). There is no further evidence indicating that the feuds in other Mudejar communities were similarly motivated, and the aljama of Castellón, in which the office of amin changed hands annually, may well have been more politically unstable. Still, given that all the communities had more or less the same political structure with officials executing the same functions, it may not be too far-fetched to infer a more general phenomenon of feuding as being in part a contest for local political power, the exercise of which could influence individual rates of taxation. Even families secure in their possession of official posts, like the Paziars, the amin s of
Alcira, were involved in feuding, which suggests that the holding of office was not the sole basis of power among the Mudejars.
Intimately linked to the competition for economic and political power as a source of feuding was a punctilious concern for honor. The Mudejar conception of honor differed somewhat from that of western Christians. For the latter, honor was attached to social rank and varied according to the possession of wealth and title. For the Mudejars, honor had the same significance at all socioeconomic levels and was the possession of the family, to be augmented or lost. This may help to explain the apparently unusual phenomenon of Mudejar shoemakers and the like fighting over points of honor.
As suggested above, family honor could be maintained and increased only through the agnates' fulfillment of their responsibilities to each other. Inasmuch as the family constituted an economic unit jockeying for its share of wealth and influence in the community, its performance in that contest reflected on its honor. However, a family without honor was by virtue of its loss of face excluded from participation in that same contest. Honor, then, was a prerequisite for the attainment of status and power in the community. Because in economic terms Mudejar society was relatively homogenous, being composed largely of small farmers and artisans and lacking an established aristocracy, the possession of honor, achieved through a family's action in accordance with the dictates of agnatic solidarity, was probably as great a determinant of local status as real wealth.
This helps to explain why so much of the violence committed by Mudejars seems to have been in defense of family or personal honor. Although economic competition might have inspired the incidents leading families or individuals to believe that their honor had been somehow sullied, it was not considered to have been in itself a sufficient cause for violence. Questions of purely economic concern were settled licitly in court; questions of honor were settled extralegally in the forum of the community. Of course, the settlement of a question of honor through violence to some extent also resolved the economic question, inasmuch as the competition was then either temporarily or permanently eliminated. Thus, in the feud between the Torralbis and the Murçis, which seems to have originated in the clash between shoemakers—Azmet Murçi against Azmet and Çahat Torralbi—Azmet Torralbi inflicted violence on Abrahim Murqi, a shieldmaker, because Abrahim had ridiculed him and thereby stained his honor. It is worth recalling that many in the morería recognized that Abrahim had incurred Azmet's wrath because of his loose tongue.
Another incident involving Muslims of Bétera shows the Mudejars' extreme sensitivity where their honor was concerned. Ubaydal Allepus
and his convert brother Miguel Crestia went to a hamlet near Bétera where they intended to mow grass with a sickle that a Christian hosteler had given them. There, a Muslim named Raboça accused the two brothers of having stolen the sickle from him. Apparently, Ubaydal and Miguel felt they had been defamed and their honor challenged, for they immediately tried to strike Raboça. Raboça then called for his brother-in-law, Amet Biari, at which point a brawl ensued that resulted in Amet's death. Although Ubaydal denied any previous acquaintance with Amet Biari, the latter's widow claimed that Ubaydal had harbored ill will against her husband. Perhaps Ubaydal's spontaneous violence in defense of his honor was the culmination of a long-standing controversy with Biari's family.
Because honor was essentially a social value, the possession of which depended on the community's evaluation of the conduct of an individual and his family, acts that entailed a challenge to or a defense of honor had a meaning recognized and understood by the entire community. Violence begot violence because social norms demanded that vengeance be exacted if honor was to be maintained. Thus, in the aforementioned stabbing of Abdalla Çentido by Açen Muça, Abdalla's friends expected that Açen would be seeking revenge. Açen, it seems, was a somewhat reluctant avenger, but he felt compelled to act when Abdalla passed by him in the street three times making insulting gestures and faces calculated to shame him publicly.
If the childrearing methods of Axa, the wife of Abdalla Murçi, are any indication of a widespread phenomenon, then Mudejar children were from an early age socialized in the ways of violent initiative and riposte in the pursuit of honor. When her nine-year-old son came home weeping after having been hit by the son of Alfona, Axa deemed that he had been shamed. She upbraided the boy for not striking back, and she demanded retaliation: "Look, when he [the son of Alfona] passes, hit him with a rock." If this were not enough, she even nagged at her husband and servant: "If you do not strike either the husband or the wife [the Alfonas] I will not consider you men. If you do not do it, go to the devil!"
A consideration of the position of women in Mudejar society confirms the continued importance of Arabo-Berber attitudes and social structures. To some extent women were pawns manipulated by their male relatives in the politics of marital alliance. Ideally, a woman was kept within the lineage through endogamous marriage, so that her lineage benefited from her reproductive power. However, if exogamous marriage were unavoidable, it was better and more honorable to receive than to give a woman in marriage. Families that had to give away their daughters in marriage did so to those families with whom alliance would
prove most useful in the local scheme of feuding relations. Although the woman in an exogamous marriage lived with her husband's lineage, she still maintained important connections with her father's family. In fact, her behavior, particularly her sexual conduct, affected the honor of her father's family, not that of her husband. This ambivalent position of the woman, bearing children for her husband's family while being responsible for the honor of her father's family, is reflected in a number of ways in the documentation.
First, because the children of a marriage belonged to the husband's family, Mudejar widows lost custody of their children, who were given into the hands of male guardians, presumably the agnates of the deceased husbands. The strategy here was to ensure that the children, along with the inheritances their fathers had bequeathed to them, continued to adhere to the father's family when their mother remarried. This explains why Fotayma resided in Sot with her new husband, Amet Albaytar, while her daughter Axa lived with two guardians in Cuart, the home of her deceased husband. Apparently, the only way a widow could continue to play a significant role in the lives of her children was either not to remarry, or to do so only to a man of the same town.
The wife's continued close ties with her father's family divided her allegiance and seem to have contributed to the instability of some exogamous marriages. This is indicated by the fact that when a Mudejar wife separated from her husband she usually returned to the home of her father. When Suçey, the wife of Abrahim Çuleymen of the morería of Valencia, went to her father's home in Petrés to attend her brother's wedding, she never returned, "her father, mother, and brother detaining her and not allowing her to come in the power of her husband." Çoltana departed with her father from the seigneury of Castell de Castells without license of the lord, even though her husband was still living there. When Ali Mançor changed vassalage from Benimuslem to Castellón de Játiva, his wife refused to accompany him and demanded the payment of her bridewealth. The bailiff of Játiva did not quite comprehend what was happening and asserted, as a Christian might, "the wife has to follow the husband wherever he would wish to go ... and to live." Probably her confidence in and attachment to her agnates allowed this woman such freedom of choice.
Although the father's household offered a haven to a woman in case of an unhappy marriage, it also harbored the harshest judges of her sexual misbehavior. An adulterous woman's shameful behavior affected mainly the honor of her father's lineage; that of the husband's remained largely unstained. Therefore, it was the responsibility of the woman's agnates to punish her. In some traditional Arabo-Berber societies the agnates were expected to kill her, and Islamic law required that adul-
terers be stoned to death. The Valencian documentation records a surprisingly large number of cases of Mudejar adultery, in which the Islamic death penalty was normally commuted by the Christian authorities to enslavement to the king. This indicates not that Muslims involved themselves in adulterous affairs any more than did Christians, but that because the honor of the father's lineage was involved in the case of an adulterous wife—and those prosecuted were almost all women—her agnates were themselves especially eager to prosecute her so as to erase their own shame. Since royal law would not permit the agnates to dispose of the woman themselves, they had to go through the proper legal channels. Unfortunately, the documents do not reveal who brought the adulteresses before the qadi 's[*] court (the criminal penalty itself had to be executed by the local Christian authority), but, all else being consistent, the agnates, not the husbands, seem the most likely candidates.
The social origins of Mudejar prostitutes further substantiate the importance of the woman's relation to her agnatic group, for the position of such women seems to have derived precisely from their having lacked the support of their agnates. Fotayma was an orphaned maidservant maltreated by her Muslim master. She ran off with a male servant who subsequently became her procurer. Mariem had left her husband, but, because her mother was forcing her to return to him, she departed from Alacuás for the brothel of Valencia. Nuzeya of Oliva had also separated from her husband, and since her parents were dead she, too, was compelled to earn a living in the brothel. Xuxa left her husband in Villamarchante for a lover who later became her procurer. Having committed adultery, she could not hope for forgiveness from her father's family. Adulterous women and prostitutes were the outcasts of Mudejar society.
If fathers were preoccupied with the sexual conduct of their married daughters, so much more were they anxious to defend the chastity of their unmarried ones. The violation of a daughter's chastity, committed with the consent of the daughter or not, constituted an assault on the family's honor. Daughters could not bring honor to a family; they could only bring shame through sexual impropriety. Thus, Aragonese Mudejars conducting business in Zaragoza brought their daughters with them and kept them secluded in the fonduk so that they would not be "maltreated" or spoiled for marriage. The Mudejars' concern to guard their daughters' chastity in defense of family honor may be explained as a reaction to an unusual phenomenon encountered in the documentation, namely, the abduction of women. In pre-Islamic Arab society, and also to some extent among Arabs in Islamic times, women were
abducted in order to shame the victimized family. The honor of the abductor's family was at the same time increased. Mudejars in royal and seigneurial morerías also seem to have employed abduction as a tactic to disgrace their enemies. Although in some cases abduction might have been, in fact, only an elopement, which still would have shamed the father who could not control his daughter, in other cases the participation of more than one man in the abduction indicates a deliberate intention to dishonor the woman's family. Mahomat and Omeymet Maixquarn of Valldigna paid a 340s settlement to the bailiff of Játiva "for having kidnapped Çayma, Mooress, daughter of the amin[ *] of the place of Manuel." The knight Miguel Çetina was sent to the Vall de Villalonga to search for the daughter ("mora donzella") of a Muslim of Millena who had been abducted by Muslims of Cocentaina. The fate of the abducted women is unknown, but it is likely that abduction lessened their prospects for a good marriage.
The passive role of women (outside of the home at least) in the maintenance or loss of family honor was linked to the economic and local political considerations underlying much of Mudejar feuding. The primary factor at stake was the woman's reproductive power. Children, particularly sons, increased a family's economic potential, and the family with many sons and strong casabiyah[*] was a force to be reckoned with in the community. Thus endogamous marriage was preferred to keep offspring within the lineage, while if exogamous marriage was necessary, it was preferable to receive the woman of another family. The widow lost control of her own children to her late husband's family for the same reasons. Adulterous unions were frowned upon because the bastard offspring did not belong to any lineage and were of no help to anyone. Prostitutes, bereft of honor and family ties, existed on the margin of Mudejar feuding society.
Above the level of family or lineage the most important solidarity among the Mudejars was that formed by the rural village or urban morería . During the era of Islamic rule such solidarity had greater force owing to the settlement of particular localities by individual clans (thus the prevalence of "Beni" in Valencian toponyms). The dissipation of larger tribe or clan solidarities and the post-Christian conquest seigneurialization of the countryside significantly modified the basis of local identity. The Mudejar no longer identified himself as the member of a particular clan, but as the vassal of a particular lord. Seigneurial control of a rural community meant that the lord's interests largely determined those of his vassals. So long as Muslim vassals stayed put, their lords tended to defend them, but not without a large degree of self-interest (see chaps. 1 and 2). The vassals themselves acted as a unit, and appear
in the documentation as the aljama of a particular place of which a particular nobleman is lord. The community of interest between vassals, and between vassals and their lords, was necessitated by economic pragmatism and a scarcity of resources. Labor, water, arable land, and livestock were all in short supply. Efficient exploitation of available resources demanded communal cooperation, and retention of these resources required collective action for purposes of communal defense. Endogamous marriage, creating large extended families, and exogamous marriage, binding different families together, both would have served to reinforce the community's unity of purpose.
One manifestation of the ability of Mudejar communities to overcome their internal differences and present a united front to their antagonists was their aggressive and joint defiance of the authorities, particularly of those sent to their villages to make arrests or to confiscate goods. When a royal constable and other officials arrived in Callosa, where they were to collect 19,666s 8d from the Muslims for pensions owed, they discovered that the Muslims had hidden their goods in places nearby. As they were returning with the goods from one of these places, the Muslims ambushed them with a barrage of stones. Officials sent to Matet to arrest two Muslims, fugitive vassals from Gaibiel, had even less luck: "there was made a great resistance by the Moors of the said place, not only breaking open the prisons where the said constable had put one arrested Moor and carrying him away, moreover, they inflicted on the said constable many blows with swords on the staff and one on the arm by which they wounded him." Three years later the amin[*] and adelantats of Matet were still refusing to cooperate in the matter. Muslim vassals of Alcocer possessing lands in the huerta of Castellón de Játiva together cleverly constructed hidden threshing floors near the river, so that they could quickly send threshed wheat downstream without paying their agricultural taxes. When the tax collector later confiscated a horse as a pledge for the unpaid taxes, two Muslims with a lance convinced him to let the horse go.
More numerous than the instances of resistance to royal officials were the clashes between vassals of various seigneuries. Neighboring villages were frequently at odds over boundaries, possession of land, distribution of irrigation water, and other such matters that affected the livelihood of their residents. Conflict of interests issued in the courts as litigation, but not infrequently took on the more sinister aspects of theft and violence. Animosity between lords also soured relations between their vassals. It is often unclear whether the actions of vassals against communities nearby were perpetrated with the consent and direction of their lord or whether the vassals took action of their own accord. For the most part one can probably presume a concurrence between lords and
vassals in such affairs, inasmuch as any gains or losses sustained by the vassals were also felt by the lord.
A brief description of the difficulties and hostility faced by the vassals of some seigneuries should provide some sense of the necessity for communal cooperation and a minimal degree of cohesion. The Muslims of Llombay's marketing of their wheat in Alcira seems to have threatened the interests of some Muslims of the town's morería . Litigation ensued between them, and the bailiff of Alcira confiscated the wheat and pack animals of Yuçeff Carroff of Llombay. A decade later eight Muslims of Llombay murdered a rival of Alcira just outside the walls of the town, which suggests a long-term dispute. Llombay's bailiff and Muslims also rustled livestock from the pastures of neighboring Carlet, for which crime the residents of Carlet planned to enter Llombay and take action. The vassals of Carlet had already been involved in a more serious dispute with Alcudia, which had resulted from the feud between their respective lords, Gaspar de Castellvi and Pere de Montagut. The latter and two henchmen killed Silim Bono of Carlet and wounded his son. Muslims and Christians of Carlet retaliated by wounding a Christian miller and killing a farmer in Alcudia. The feud was finally resolved when Montagut married Castellvi's daughter. Some conflicts were more one-sided. The attack of the brothers Ferrer with their squires and Muslim vassals so terrorized the Muslim residents of Faldeta that they deserted the place.
Perhaps the most frequent cause of strife was a community's misappropriation of irrigation water, which placed in jeopardy the crops of other communities in the vicinity. The lord of Alginet complained that his village was "perishing" because the officials, including the amin s[*] , of the Foya de Llombay were not allowing the water to flow as accustomed. The consequences of such disputes could prove dire, such as the one between Antella and Sumacárcel over "a bridge or a duct of an irrigation canal," which provoked Muslim and Christian vassals of Antella to kill Çahat Torraboni of the rival village.
At the local level Mudejar society displays two conflicting tendencies: on one hand, fragmentation, as manifested in the feuding relations between rival family or lineage groups; and, on the other hand, a reflexive solidarity necessitated by the struggle between communities over the possession of scarce natural resources. The former tendency was rooted in traditional Arabo-Berber modes of social organization. The sources of the latter tendency are more complex, for the importance of seigneurial rivalry and the role of Christian vassals in the strife between communities disallows a description of local solidarity as a purely Mudejar phenomenon. Nevertheless, it is evident that the Muslims inhabiting individual villages were able to unite when the livelihood of the entire
community was threatened, even if at the behest of their lord and in conjunction with their Christian neighbors. That such was the case perhaps can be explained in part by the ability of Arabo-Berber segmentary societies to amalgamate when necessary, as well as to atomize.
Considerable space has been devoted to a discussion of Mudejar social structure and to an analysis of Mudejar feuding relations not only to provide some insight into life within the morería below the surface of the Muslim-Christian interface but also to make a point crucial to the understanding of the remarkable tenacity with which the Mudejars and, later, the Moriscos adhered to Islam despite considerable pressures, both informal and formal, to the contrary. The point is that for the Mudejars religious conversion involved much more than a change in their religious beliefs, a change radical enough in itself; it demanded a fundamental alteration of their social attitudes and social organization. The Mudejars were shielded from the allure of Christianity not only by their profession of a faith as exclusive as Christianity but also by the vitality and structure of their subsociety, which was founded on social practices and assumptions distinct from those of their Christian neighbors. While Mudejar feuding in its outward bloody manifestations did not differ from Christian feuding, the social significance Muslims ascribed to it and the family structures and system of values on which it was based were distinct. The feud as a process of status determination was group-specific, functioning in its very intensity to reinforce traditional attitudes and structures. Moreover, religious belief and social practice were largely coterminous, so that the Mudejars' distinct social customs were as much a sign of their "Moorishness" as was their belief in the oneness of God and the prophethood of Muhammad. For instance, conversion to Christianity would have meant a prohibition of the practice of endogamous marriage (between first cousins), a custom sanctioned by the Prophet and essential for the maintenance of their Arabo-Berber social structures. The frequent and often friendly meeting of Muslim and Christian on the neutral ground of marketplace or tavern no more resulted in a merging of their different forms of social organization than it did in religious syncretism. Indeed, Mudejar social behavior, particularly the marked propensity for feuding—or at least the style and dynamic of that feuding—was sufficiently different so as to evoke comment from among the Christians. Francesch Centelles, a shoemaker well acquainted with the Muslims of Valencia's morería , when asked to testify in court about the character of Abrahim Murçi, responded in a manner that suggests that the feuding Mudejar had become a stereotype: "he is a man who seeks fights and quarrels, like any other Moor."
It seems clear that the Mudejars' distinctively Arabo-Berber mode of social organization helped to shore up their cultural boundaries against the acculturative attrition of an overwhelming Christian presence. These boundaries preserved the essential element of their ethnic identity, the profession of Islam. Although social mores and behavior were intimately bound up with and were to a large extent the product of religious belief, nevertheless, they in themselves were insufficient to perpetuate the religious faith of the social group. The decision of Mudejars to flee the kingdom in 1502–1503 instead of abandoning Islam was the expression of an intensity of faith that transcended the more amorphous "Moorish" cultural identity engendered by repetitive social practice. It is necessary, therefore, to comprehend how the Mudejars actively instilled and fostered their Islamic faith and identity.
An essential element buttressing the faith of individual Muslims was the sense of belonging to a larger community of believers, the ummah . Whatever the situation of the ummah 's individual component polities, even those long since subjugated to Christian powers, membership in the ummah served to distinguish them from all non-Muslims. Unfortunately, the Muslims' adherence to a common faith did not preclude divisiveness within the ummah . Since the fall of the Umayyads in the eighth century the Islamic world had been rent by factionalism, and, as has been seen, this was no less true of the Mudejars. However, making allowances for human imperfection and the not unusual inconsistency between religious precept and social practice, that the Mudejars frequently embroiled themselves in feuds, despite the fact that Muhammad had inveighed against such fratricidal strife among Muslims, does not mean that they had lost their sense of Muslim identity, particularly their collective identity vis-à-vis the Christian world. On the contrary, there were a number of instances in which the Mudejars appear to have acted as a collectivity or were perceived by the Christian authorities to have been such. Let us recall how the nobility advised Fernando against the forced conversion of Valencia's Muslims (1502), noting that "they have their communications with each other," and that any untoward royal initiatives would provoke a violent mass Mudejar reaction. The Mudejars also seem to have been united in a common concern for the embattled sultan of Granada, as was manifested in their taking up collections on his behalf, praying in their mosques for his victory, and negotiating with the Ottoman Turks to come to his rescue. Likewise, the aid provided to runaway slaves by the Mudejars—not just by one community but by any number of morerías in which the fugitives hid on
their way to freedom—demonstrates their ability to act together as Muslims for the benefit of other Muslims.
Mudejar group-consciousness may be seen as the sum of each Muslim's perception of the fundamental difference between himself and his Christian neighbor, and of each Muslim's participation in the life of an autonomous community juridically framed by the Shariah, a corpus of law at once religious and secular. This aggregate awareness of individual Mudejars, however, seems in itself insufficient to have counterbalanced the animosity between feuding families and competing communities, or to have allowed for the Mudejars' alleged network of "communications" and their ability to act as almost a single political entity. One must seek a more concrete pattern of relations transversing the cleavages between agnatic lineages and rival neighboring communities.
The role of exogamous marriage in binding families together and the factor of intercommunal strife in promoting solidarity among the inhabitants of any one village has been discussed. The apparent impediment of intercommunal conflict to a larger, kingdom-wide Muslim solidarity presents a greater problem. Mudejar economic activity, establishing contacts between Muslims of all walks of life from a variety of localities, would have been a key factor in circumventing, or at least lessening, the rivalry between communities. It must be emphasized that while the Mudejars were vassals of particular royal or seigneurial morerías , their economic activities were not circumscribed by the boundaries of any one place. This was especially true of itinerant retail merchants, whose vending took them from the northern to the southern reaches of the kingdom. Since by the nature of their work they were more or less unattached to any particular local interest, these merchants would have served as an appropriate medium for relaying information from one community to another, tying together the separate worlds of distinct aljamas.
Another group whose activities were equally unhampered by specifically local concerns were the licensed mendicants, who traveled throughout the kingdom begging alms from their Muslim brethren. Since charity was one of the Five Pillars of Islam, these itinerant mendicants provided an opportunity for pious Muslims to express their religiosity in a manner unrelated to the secular aims of the family and community. Both the giver and the receiver of alms participated in a transaction that emphasized exclusively their obligations as members of an Islamic community, not those stemming from kinship or from residence in a particular place.
The economic interplay between town and countryside made the kingdom's urban centers sites for the meeting and mingling of Muslims from various rural villages. Mudejars traveled to town to market pro-
duce, to purchase the manufactures of local artisans, or to pass their leisure time in the taverns or fonduk. The large Christian populations of the towns would have induced Muslims from out-of-town to congregate with their coreligionists before turning to Christians for comradeship. If towns like Játiva, Alcira, Castellón de la Plana, and Villarreal attracted Muslims from surrounding hamlets, the capital city gathered in Muslims from all over the kingdom—indeed, from all over the peninsula. Valencia, a veritable teeming metropolis, served as a "melting-pot" for the kingdom's Mudejars. This is indicated in the documentation by the frequent appearance of non-local Muslims working, purchasing, and pursuing litigation in the capital. Furthermore, new vassals in Valencia's morería originated from seigneuries of a wide geographic range, whereas those swearing vassalage to the king in towns such as Játiva or Alcira came from seigneuries nearby. Change of vassalage in itself accentuated two other trends that created links between Muslims of different communities. First, many Muslims who changed vassalage still continued to hold and cultivate land in their former seigneuries, either themselves or through local sharecroppers. This further complicated the already complex Valencian pattern of land tenure, which saw farmers and artisans renting small parcels of land in a variety of places, not just in their place of residence. Mudejars with economic interests in diverse localities likely would have had friends and contacts of equally diverse origins. Thus land tenure itself sometimes cut across the lines of economic competition contingent upon strictly local affiliations. Second, change of vassalage caused the fragmentation of local lineage groups as nuclear families left their agnates behind when they settled elsewhere. Given the importance of c asabiyah[*] in Mudejar society, agnates living in different places were probably still bound by kinship. In addition to these intercommunal agnatic links, there also occurred exogamous marriages between families of different towns, such as the one uniting the niece of Játiva's qadi[*] to the son of Valencia's qadi . Therefore, kinship, both agnatic and affinal, created a network of interests that would have mitigated the intensity of the rivalry between communities for economic reasons alone.
The individual Mudejar's sense of belonging to a kingdom-wide Muslim community and the ritual expression of his commitment to Islam coalesced in the act of pilgrimage to the mosque of Atzeneta in the Vail de Guadalest. This mosque housed the sepulcher of the Sufi mystic Abu[*] Ahmad[*] Jac far b. Sid-bono[*] al-Khuzac i[*] (d. 1227). From the thirteenth century until 1570, when King Felipe II had the mosque of Atzeneta destroyed, the tomb of this saint attracted Muslim pilgrims from all over the kingdom of Valencia, and sometimes from Aragon, Catalonia, Granada, and the Maghrib as well. Ecclesiastical views, expressed at the
Council of Vienne in 1311, that such Muslim pilgrimages were an affront to the Christian community, combined with royal misgivings about large numbers of Muslims of diverse origins gathering each year at the shrine of Atzeneta by 1379 resulted in Crown attempts to prohibit this pilgrimage. The royal authorities, however, were unsuccessful, for throughout the fifteenth century and much of the sixteenth century Valencian Muslims continued in "semiclandestine" fashion to journey to the tomb of Sid-bono[*] . This annual act of Islamic devotion thus became for the individual Mudejar a statement of resistance to Christian authority, a politically dangerous affirmation of identity with the other participants in the pilgrimage.
As evinced by the Mudejars' awareness of the events occurring in the wider Islamic world, especially in Granada and the Maghrib, and by their political activities in conjunction with Granadan, Maghriban, and Ottoman Turkish Muslims, their understanding of what constituted the community of believers extended far beyond the borders of the kingdom. Recognition of this much larger community alleviated their sense of isolation and hopelessness, and strengthened their own Muslim identity, particularly when the large majority of that community was governed by Muslim rulers, some of whom, like the Ottoman sultan, were extremely powerful. As was pointed out in chapter 2, concrete family and commercial connections underlay the Mudejars' politicoreligious identification with the dar[*] al-Islam[*] . As a consequence of previous Mudejar emigration to Granada and the Maghrib, Valencian Muslims had family branches in Islamic lands. On account of these kinship ties, Mudejars traveled to Almería, Tunis, and Oran for the purpose of collecting the inheritances left them by deceased relatives, or, like Ali Fotoffa of Bétera, in order to visit those still alive. A Muslim family of Cocentaina journeyed to Granada to attend a family wedding, while a widow of the morería of Valencia married a Granadan Muslim and then departed with him to North Africa. Like this marriage, there were other instances of recent emigration that forged new links of kinship between Valencia and the Maghrib. The sub-qadi[*] of Játiva decided to spend his retirement in the Maghrib, while his son stayed behind in Játiva and succeeded him in office. Yahye Bellvis, the brother of the qadi general, moved to Tunis and continued to benefit from his commercial connections in Valencia. It was precisely such ties of kinship that facilitated Fernando's settlement of Granadan Muslims in Valencia after the conquest. Thus, Mahomat Fuçey of Bellreguart was licensed to travel to Almería "in order to fetch some relations that he has in the city." Reciprocal commercial interests strengthened Mudejar affinity for Granada and the Maghrib. Mudejars journeyed to Almería and Tunis to sell their merchandise, while Maghriban mer-
chants came to Valencia on business. The royal licenses that permitted these merchants to reside in the kingdom for a year or more created ample opportunity for contact with Mudejars. Religion, kinship, and commerce all bound the Mudejar inextricably to the dar[*]al-Islam[*] . It is doubtful that the Mudejars' Muslim identity and group consciousness would have fared as well had they been isolated.
Relative isolation of a different sort actually abetted the Mudejars in their preservation of an Islamic culture. Historians have puzzled over the fact that despite royal efforts to attract Muslims to urban royal morerías , where the tax burden was lighter, the large majority nevertheless preferred to remain on seigneurial lands. The Mudejars' choice of residence is best interpreted as having had a religiocultural foundation, rather than an economic one. Life in the largely Christian cities posed obvious threats to the integrity of Mudejar Islamic culture. Either militant Catholics were endeavoring to eradicate all signs of Islam—calling for the destruction of minarets, prohibiting the call to prayer, and the like—or the pleasantries of city life were insidiously weakening the Muslims' resolve to live in accordance with the Shariah. That pious Muslims were sensitive to the latter threat is evinced in the complaint of the aljama of Játiva regarding the nocturnal activities of Christian youths in the morería and the deleterious effect that the "dishonest dress" of the alfondeguer 's wife might be having on Muslim youths.
In contrast, life in seigneurial villages afforded the Mudejars a refuge from an aggressive and expanding Christian presence. In these villages Muslims sometimes composed the majority of the population, and their freedom and comfort in religious observance were correspondingly greater. The lords seem to have had few qualms about the public manifestations of Islamic worship. For instance, they allowed their Muslim vassals to make the call to prayer with a horn and perhaps vocally, whereas it was prohibited in the cities. During the time of the Moriscos, the seigneurs were infamous for permitting their ostensibly Christian vassals to practice Islam and for protecting them from the Inquisition. The Mudejars also seem to have benefited from their lords' religious tolerance, even if at a price. For most Mudejars the religious freedom thus secured was sufficient compensation for the heavier burden of seigneurial dues.
Not surprisingly, the centers of Islamic learning in Valencia, such as they existed, were for the most part located in rural villages, not in urban morerías . Of the twenty-five Mudejars who journeyed to Almería, Tunis, Oran, and Granada for the purpose of studying the Arabic language and Islamic law (see table 19), only five came from urban morerías —four from Játiva and one from Castellón. The others all came from seigneurial lands—Ondara, Cuartell, Artana, Mascarell, Valldig-
na, Benilloba, the Vall de Uxó, and so on. Barceló Torres, intimating the cultural inferiority of the urban morerías , notes that of the 270 Mudejar and Morisco Arabic documents she has found, only nine were drawn up in the morería of Valencia. It is indicative of this state of affairs that Çahat Coret of the Foya de Buñol, who "applied himself diligently to Agarene letters," was appointed faqih[*] of Valencia, after the aljama had failed to find anyone in the morería sufficiently learned to fill the post.
It appears, then, that the Mudejars' Muslim identity was nurtured both through the unobstructed public worship of Islam, a freedom they
secured by their choice of residence, and through the maintenance of and participation in a literate Arabic culture. Clearly, the latter was needed to sustain the former. Even a minimal level of popular religious awareness necessitated the mediation of learned men (c ulama[*] ') who could read and interpret the Qur'an[*] for the faithful (the Arabic dialect spoken by Mudejars was different from the classical language of the Qur'an). Beyond that, men conversant in jurisprudence (fiqh ) and all that entailed—a knowledge of the Qur'an and of the customs of Muhammad and his companions (Sunnah ) as set down in the traditions (hadith[*] )—were needed to administer justice in the Islamic courts, either as qadi[*] or as faqih[*] , and to see to it that the community lived as much as possible within the framework of the Shariah. Taking into account the Mudejars' situation as a minority enclave composed primarily of farmers and artisans, the grooming of even a small group of culama[*] ' required a determined and sustained effort. Mudejar acquisition of the necessary cultic and legal knowledge was in itself a considerable achievement. The social and intellectual environment was unpropitious for the creation of original scholarly works.
Arabic instruction given to Mudejar children in local schools perpetuated this rudimentary but essential Arabic culture. The thirteenth-century capitulations had granted the basic privilege of maintaining schools to the Mudejars. The Muslims inhabiting the new morería of Orihuela (formed in 1446, but lasting only until 1451) were allowed "to have a schoolmaster." The morería of Valencia also had a school, at least until 1455. Documentation from Fernando's reign contains references to schools operating in Ondara, Oliva, and Valldigna. Perhaps Çahat Coret of the Foya de Buñol began his studies of "Agarene letters" in his hometown. Since the Mudejars were able to give at least an elementary Arabic education to their children in Valencia, one may infer that the Mudejars who took the trouble to travel to Granada or North Africa for study did so not merely "to learn to read and write Moorish," as the travel licenses state, but to pursue more advanced studies, particularly in jurisprudence. Tunis and Almeria were both well equipped to meet the academic needs of the aspiring Mudejar faqih[*] .
It is difficult to know whether there were schools for more advanced studies in Valencia, although it seems that the Mudejars possessed a sufficient amount of learned Arabic works to have allowed for at least the informal meeting of erudite culama[*] ' with students eager to learn. Juan Andres, the convert from Játiva, recounted that his father, a faqih[*] , had taught him jurisprudence. Barceló Torres's search for the bits of Arabic literature surviving in Valencia reveals that fifteenth-century Mudejars had access to Qur'ans, hadith literature, devotional works,
and legal works. Also, in 1450 a faqih[*] of Paterna brought back from Cairo a treatise on trigonometry, in which the use of an astronomical instrument is explained. The most impressive information on Mudejar higher learning comes, surprisingly, from the kingdom of Aragon. The letter of a student to a faqih in Belchite reveals the existence of a madrasah (school) in Zaragoza as late as 1494. There the student in question studied theology, and medicine from the Qanun[*] of Ibn Sina[*] (Avicenna). Considering that such a school still functioned in Aragon, where the Muslims' fluency in Arabic was much less than that of their brethren in Valencia—although perhaps not as minimal as was once thought —it would seem that similar centers of advanced study must have existed in Valencia as well. It is doubtful that every Valencian faqih had the opportunity to travel to Islamic lands for study; some were probably purely local products. Moreover, there were Mudejar physicians and surgeons in Valencia, and these professions required a certain amount of learning, perhaps in the classical Arabic medical texts. Juçeff Alatar, a surgeon of Valencia, was granted a royal license to practice after administering to a Christian knight and passing the examination given by a Christian "master in medicine."
The Mudejars frequently utilized their Arabic literacy in a far more mundane fashion in the writing of letters and contracts for official and private business. The extant Arabic documentation contains records of tax payment, and letters to and from local amin s[*] concerning the collection of taxes and debts from Muslim vassals or judicial procedure against them. Much of this correspondence was between amin s and royal bailiffs, which indicates the functioning and interpenetration of two levels of bureaucracy: the all-encompassing royal Romance-Latin administration and the local Mudejar Arabic administration manned by amin s, qadi s[*] , and faqih s. It has already been demonstrated how the two bureaucracies interrelated in the matter of the Crown's taxation of Mudejar inheritances (chap. 4). The Christian authorities' recognition of Arabic documents as valid evidence in litigation and as contractually binding in business transactions, even those between Muslims and Christians, gave Arabic an "official" status in the kingdom of Valencia. For instance, when passing sentence in favor of Fatima Bisquey, who claimed that she owned half of a house given to her as bridewealth (sadaq[*] ) by her husband, against the opposing claimant, the merchant Berthomeu Pinos, the bailiff general pointed out that the decisive evidence was "an act of acidach (sadaq[*] ) and/or marriage contract—exhibited in the trial on behalf of the said Fatima—received by the qadi[*] and/or faqih[*] Mahomat ben Abdulaziz Alcari on the date of 11 March 892 of the Moorish calendar." Muslims bound themselves to pay debts to Christian creditors by acknowledging their debts in Arabic
documents. In an Arabic document written in his own hand, Ubaydal Donzell confessed, "I, Ubaydal Donzell, recognize that I owe to you, Manuel Bou, eighteen and one-half pounds, which are for spices and alum."
The cultivation of Arabic for higher intellectual pursuits—Qur'anic[*] study, fiqh , medicine, and so on—for the drawing up of legal instruments of various sorts—marriage contracts, letters of debt, and tax records—and for daily parlance lent the Mudejars a common ethnic identity and group consciousness on the basis of language alone. Their knowledge of Arabic allowed them to participate in the intellectual life of the wider Islamic world, just as their understanding of Romance enabled them to function more efficiently in Valencia's Christian society. The Mudejars' use of Romance, however, was far more occasional, employed only when they desired or needed to communicate with Christians. Otherwise, Arabic was an effective social and intellectual barrier between Valencia's Muslim and Christian communities. Indeed, the Mudejars' use of Arabic sometimes aroused Christian suspicions. When Muslim slaves escaped from their masters, any Mudejar who had been seen speaking with the slave in Arabic was considered a prime suspect as an accessory to the crime. The Christians assumed that the Mudejars' choice of language defined their sympathies and guided their actions as persuasively as did their religious faith. There was much truth in this assumption. Since Arabic was the language of the Qur'an[*] , literally the word of God dictated to Muhammad, its use by the Mudejars had a special spiritual significance, and therefore contributed to their perception of themselves as Muslims. The veneration of the Arabic language itself explains the Aragonese and Castilian Mudejars' and Moriscos' writing of aljamiado literature (Romance written in Arabic script) as a means of strengthening their Muslim identity. It also explains why the Christian authorities decided to prohibit the Valencian Moriscos' teaching of Arabic to their children as a means of effecting their true conversion to Christianity.
Social structure, language, communal and judicial autonomy, ties of kinship and commerce within Valencia and with their coreligionists in the dar[*]al-Islam[*] —all contributed to the Mudejars' distinct religioethnic identity and to their perception of themselves as a single body united in stark cultural opposition to Christian society. Still, the body required animation and direction, a sense of purpose particularly Islamic. This was provided by the faqih s[*] . They functioned in the Mudejar social body as spiritual cadres, infusing its individual communal cells with a commitment to Islam and, by virtue of their grounding in a common intellectual tradition and world view, binding those cells together in a unity of religious purpose. As the local fonts of religious and legal knowledge, the
jurists were eminently suited for this task. By offering their legal opinions and resolving disputes on the basis of the Shariah, they ensured that it remained the lofty standard against which Mudejars evaluated their own conduct and by which they endeavored to regulate their lives.
The plea of the newly converted Moriscos in their negotiations with Carlos I in 1526 reflects the great esteem in which the faqih s[*] were held by the kingdom's Muslim populace. The Moriscos informed the king that in the days before the conversions, "when the call to prayer was made in the mosques," the Muslims of the kingdom used the rents from those properties bequeathed by the pious to the mosques to pay the salary of the faqih s, "who have consumed their whole life in studying and knowing the Moorish law and have not been concerned with other offices." The Moriscos went on to request that a portion of the rents pertaining to the new Morisco churches, formerly mosques, continue to be reserved for the support of the baptized jurists. The Morisco faqihs —and they are still referred to as alfaqins in the sixteenth-century documentation—continued throughout the sixteenth century to form the core of Morisco resistance to the official Christian program of religious and social assimilation.
The aforementioned Morisco request reveals important information about the Mudejar jurists. It is clear that they devoted their entire lives to the study of Islamic law and, presumably, of its foundations, the Qur'an[*] and the Sunnah . Because the study of the jurists ensured the continuity of Islam as a living religious and intellectual tradition, the Mudejars deemed it an essential activity, so essential that they used the pious endowments (waqf ) bequeathed by the faithful to the mosques to support the jurists. Furthermore, the jurists were supported in such a way that they would not need to bother themselves with any labor other than that properly religious and legal. It is important to note that in Islam there were neither priests nor an ecclesiastical hierarchy for whom financial support was institutionalized, as was the case with the Catholic Church. The Mudejars' support of the faqih s was made possible by the will of the community, a local adaptation to a situation in which Islam had long been deprived of public primacy.
Despite the faqih 's key role in the life of the Mudejar community, the documentation provides, unfortunately, very little information about him. The reason for this is that the faqih 's sphere of activity—the mosque, the madrasah , the Islamic court—was very rarely impinged on by the Christian world. Matters concerning Muslims alone and not affecting the public life of the kingdom had no importance for the king and his officials. As a result, it is the amin[*] , the fiscal and juridical intermediary between Muslim and Christian worlds, who appears most often in the
documentation. The activities of the faqih[*] held as little interest for the Christians as did the Muslims' theological views.
The primary role of the faqih was that of jurisconsult, acting either as counsellor or as arbitrator in litigations between Muslims. Ageg b. Çaat Ageg of Alcira paid 10s to the faqih of Villalonga for having counseled him in his dispute with Abraym Xativi. A faqih of Ondara traveled to Ribarroja in order to "treat with and reconcile a Moor and a Mooress, husband and wife, who wanted to separate." The faqih also taught in the local school, a task for which his years of study had well prepared him.
Also, as the Moriscos stated in 1526, the faqih s "were serving in the mosques," undoubtedly as preachers. Regarding the content of their sermons, the only information comes indirectly from the allusions made by the nobility when, in 1502, they beseeched Fernando not to convert their Muslim vassals. It seems reasonable to assert that the sermons of the jurists comprised primarily instruction in the basic tenets and precepts of Islam and positive exhortation to conduct oneself accordingly. As a consequence of their teaching, the lords stated, "among them [the Mudejars] each one defends the said sect [of Islam] and has worked and works [to the end] that the Moor may be a good Moor." More interesting is a type of sermon more negative in tone, which constituted a defense of Islam through a disparagement of Christianity. The nobles offered an illuminating explanation of why the Mudejars were "beside themselves" with fear that the Inquisition would proceed against them all:
they [the Mudejars] say that ... none of them could be excused [from prosecution by the Inquisition] because publicly they have had and have in the present kingdom their mosques and their faqih s, who publicly admonish them that the sect of Mahomat is better than the law of the Christians and that all [Christians] end in this damnation.
By emphasizing the threat of eternal damnation, the jurists hoped to discourage potential Muslim apostates who might be toying with the idea of baptism for worldly reasons. This kind of preaching may be interpreted as a reaction to the pressures exerted by an increasingly militant Spanish Church and to the threatening pervasiveness of Christian culture. The tragic end of the Spanish Jews and Conversos, and the incipient movements of the Inquisition against Islam, could hardly have failed to impress upon the faqih s the necessity of a defensive anti-Christian posture.
The contemporary struggle between Christian and Islamic states did
not fail to influence the direction of the faqih s[*] ' activism. The jurists understood what implications the conquest of Granada might have for the future of Islam in the Iberian peninsula, in terms of both the morale of the Mudejar populace and the Mudejars' treatment by the Catholic Monarchs. Consequently, the jurists' activity, particularly their preaching, took on a markedly political tone. It was they who collected in the morerías funds for the aid of the beleaguered Nasrid sultan. Moreover, the king was informed:
the said Moors and the faqih s of the said morerías since the time of this enterprise [the war against Granada] have ordered a certain prayer and they make that [prayer] continually in their hours [of prayer], [the prayer] containing, in effect, that God should exalt the said king of Granada and that He should destroy us [King Fernando] and all our people.
While the efforts of Mudejar jurists could hardly have altered the course of political events, they nevertheless succeeded admirably in strengthening the commitment of their congregations to Islam. This commitment is evinced in the very low rate of Mudejar apostasy. One document offers a rare glimpse of Mudejar sensitivity in matters of faith, an area of life where Christian interference was not easily endured. When royal officials made the mistake of entering a mosque in Ondara in order to apprehend a condemned Muslim criminal, the reaction of the congregation was one of violent indignation. "The Moors, amin s[*] , jurates [adelantats ], and all the people" hurled "stones and ... the tiles from the roofing, and with lances and crossbows wishing and working to damage you [the lieutenant governor] and your ministers, made a great resistance against you, perturbing you and preventing the capture of the one convicted."
The success of the faqih s in overcoming intracommunal and intercommunal factionalism and imparting to each community a sense of commitment to the common cause of Islam was furthered by two factors: the communication between those possessed of religious and legal knowledge, and the foundation of the jurists' status on terms different from those which determined the prestige of other community members. Even though each faqih belonged to a particular community, the faqih s did not exist in a state of intellectual isolation, instructing their congregations and offering legal opinions without consideration of the opinions and perhaps greater knowledge of their learned fellows in other places. On the contrary, their role as the transmitters of a common tradition and their very similar intellectual formation—all having been educated in Granada, the Maghrib, and Valencia—facilitated consultation among the learned and, indeed, advised it, if they were to
maintain a consistent orthodox standard. The activities of the faqih[*] Abdalla, originally a captive from Tripoli, suggest a network of communication and consultation among the Muslim judges and jurists of the kingdom. While in Valencia, Abdalla met the faqih of Manises, and they discussed Islamic law. The latter then invited him to dinner in Manises. He was also a friend of the qadi[*] of Benaguacil and lodged in his home. In Ribarroja Abdalla acted as marriage counselor to an unhappy couple. He taught school in Ondara, Oliva, and other places, and while in Oliva he conferred with the faqih and "read in the said morería ." More interesting still, he and the faqih of Paterna exchanged Arabic books. Abdalla's career was probably somewhat more peripatetic than that of most jurists, since he had to wander about collecting alms to repay the aljama of Ondara for having ransomed him. Nevertheless, Abdalla was able to "go among the faqih s of the present kingdom begging for the love of God," because there were established channels of communication among the learned, and because he himself was "a man of knowledge."
The career of Abdalla also demonstrates that the Mudejars on the whole respected and heeded the opinions of learned and holy men. Abdalla was known by Muslims throughout the kingdom and was reputed to lead the life of a saint. The qadi of Játiva, the kingdom's largest morería , related why he and the aljama wished to make Abdalla their jurist. Abdalla, the qadi pointed out, "is a very good Moor and ... leader of prayers [oracioner ]," so much so that after the death of the former faqih of Játiva, Abdalla, owing to his "good fame, life, and knowledge," was the unanimous choice to succeed him.
Clearly, the prestige of Abdalla, a foreigner and technically a slave, and of men like him in the eyes of the Mudejar community, rested neither on wealth nor on family backing; rather, their status and influence, both local and, in the case of Abdalla, kingdom-wide, derived from their knowledge of religious and legal tradition and from the holiness of their lives. This is suggested in two other cases mentioned above: the intervention of a faqih from Villalonga as legal counselor in a litigation between Muslims in distant Alcira, and the appointment of a Muslim of the Foya de Buñol as faqih of Valencia on the basis of his diligent studies alone. Put in another way, the jurists were able to rise above the petty feuding between Mudejar families and villages because their way of life, financially supported by the community, removed them from the competition for honor, wealth, and political power. Their opinion was heeded because they were nonpartisan and had no stakes in that competition. It is probable that the jurists were able to attenuate the intensity of local feuding by acting as legal counselors and arbitrators between disputants. Their preaching was persuasive because it was an expression
of the knowledge and piety that so few possessed. The jurists were able to appeal to the Mudejars as Muslims on a level of consciousness unrelated to mundane local concerns. Because the jurists themselves had a scholarly network of sorts and a certain consensus of opinion regarding the Mudejar community's needs, they preached a similar message. Their message, that of Muslim resistance to Christian assimilative pressure, had efficacy and resonance because, on one hand, they themselves were dispersed throughout the kingdom's Muslim population, local products closely tied to their people, while on the other hand, they possessed vital knowledge that raised them above the mass of Mudejars in an overarching network of religious leadership.
For reasons beyond the control of Valencia's faqih s[*] , the Muslims of the kingdom were to spend their final years in the peninsula as unwilling Christians. As has been seen, the chain of events that led to the Mudejars' conversion began with the fall of Granada. This signal and seemingly conclusive event in the long and bloody peninsular struggle between Islam and Christianity had an unforeseen and somewhat ironic consequence. For the Mudejars of Valencia it resulted in a cultural windfall that would help them in the difficult days ahead.
When Fernando settled conquered Granadan Muslims in Valencia and promoted the sale of Maghriban and Granadan prisoners of war within the kingdom, he was not only helping himself by increasing royal revenues but also unconsciously contributing to the Mudejars' Muslim identity and group solidarity. The influx of numerous Muslim captives—hundreds from Málaga alone—elicited from the Mudejars considerable cooperative effort on behalf of their Muslim fellows. Mudejar communities collectively ransomed Muslims and, in a seemingly organized fashion, provided assistance to runaway slaves. More important still was the type of Muslim brought into the kingdom through Christian conquest and piracy. These Muslim captives and settlers, unaccustomed to Christian rule and little affected by the acculturative impact of long-term coexistence with Christians, were in all likelihood more steadfast in their commitment to their ancestral faith. Among the new arrivals were men of learning. There were physicians from Granada, jurists and readers of the Qur'an[*] from both Granada and the Maghrib, and a Sufi mystic from the Maghrib. True, it is not known what became of these men, although the career of the faqih Abdalla of Tripoli, a captive known throughout the kingdom for his learning and piety, is, if not typical, suggestive. In this regard, a puzzling but interesting comment was made by some Christians about Abdalla. When the Christian hostess of a hostel and some Muslims lodging there introduced Abdalla to some Christian guests as a faqih , the Christians remarked,
apparently in jest, "he is black and he could be a faqih[*] ." Given the similarity between Iberian Muslims and Christians in terms of skin color, the reference to Abdalla's darker coloration as if it were part of a widespread stereotype of Valencia's faqih s hints at a perhaps more general phenomenon of Mudejar religious leaders with Maghriban origins. In any case, the social organization of Mudejar society would have facilitated the integration of Granadan and Maghriban Muslims, both erudite men and less extraordinary folk. Thus, Mudejar society was strangely reinvigorated as a result of the Monarchs' war against Islam. While the conquest of Granada ushered in the tragedy of the Moriscos, it also ensured that the Christian authorities of Valencia would have more formidable opponents in their struggle to eradicate Islam from the kingdom. Islam survived in Valencia not as a fossilized remnant of thirteenth-century Almohad glory, but as a resilient and adaptive society, steeled by its social structure and inspired by its faqih s. |
Easy Magic Tricks for Scouts
> > > Triple Prizes < < <
piece of paper
- Write 1089 on the paper without showing anyone, fold it, and place it on the table in plain view.
- Give someone a piece of paper and pencil.
- Tell them to write down any 3 digit number that uses 3 different numerals in the middle of the paper. Not 111 or 202 or 330 where the same numeral is used more than once.
- Tell him to reverse the number. If the number is larger, write it above the first one. If smaller, write it below.
- Subtract the smaller from the larger.
- If the resulting number has 2 digits, fill in ahead of it with a zero.
- Reverse the number and write it below the bottom number.
- Add the bottom two numbers.
- Unfold your paper and ask if it matches their result - 1089!
Guessing the Coins
- Give the penny and nickel to a friend.
- Turn your back and tell him to hold one coin in each hand.
- Say that you want him to focus on the coins and mathematics is the best way to get his mind in gear.
- Ask him to multiply the coin in the right hand by 2 and to say 'OK' when he has the answer.
- Ask him to multiply the coin in the left hand by 17 and to say 'OK' when he has the answer. If it took about the same amount of time for each, then the nickel is in the right hand since 1x17 is easy to compute. If the 2nd one was much slower, then the nickel is in the left hand.
The next time you do the trick, use different numbers like 13 or 19.
- Give your friend these instructions:
- Think of any number from one to 10.
- Multiply it by 9. [Pause while they do this]
- If it's a two-digit number, add them together.
- Now, subtract 5 from the number in your head. [Pause again]
- Now, think of the letter in the alphabet that corresponds with the number you are thinking about. For instance, if you are thinking of the number 1, it would be "A". Number 2 would be "B". 3 is "C", and so on.
- Now, think of a country that starts with the letter you're thinking of.
- Spell the country in your head. [Pause here]
- Think about the second letter in that country's name. Now, quickly think of an animal who's name begins with that letter.[Pause here]
- Now, think of the animal's color.
- [Pause and concentrate] That's funny... this can't be right... there ARE no gray elephants in Denmark!
Add to 15
paper and pencil
- Write your answer (15) on a slip of paper, fold it, and lay it on the table.
- Give your fiend the paper and pencil.
- Ask him to draw a Tic-Tac-Toe board.
- Ask him to write down the numerals 1 to 9 in order in the diagram. 1,2,3 on top. 4,5,6 middle. 7,8,9 on bottom.
- Have him circle any numeral in the first row.
- Have him circle any numberal in the second row that is not in the same column as the first one circled.
- Have him circle the numeral in the bottom row that is in the column with no circled numerals.
- Have him add up the circled numbers.
- Have him unfold your answer - 15!
Jump the Coin
- Bet your friend that he will not be able to jump over this coin when you set it down.
- Have him stand with his arms stright down at his sides and then place the coin on top of his head.
- If he complains that you tricked him, then give him another chance - put the coin in the very corner of the room.
Guess the Color
4 pieces of paper - 3 inch square
red, blue, and green markers
- Put a red spot on one piece of paper, blue on another, green the third.
- On the back of the red paper, write in small neat printing - "YOU CHOSE RED"
- On the envelope where the stamp would go, write "YOU CHOSE BLUE" - not too big so you can cover it with your thumb.
- On the 4th piece of paper, write "YOU CHOSE GREEN", fold it, and put it in the envelope.
- Put the 3 colored papers in the envelope and now the trick is ready.
- Pull the envelope out of your bag covering the writing with your fingers and lay it front down on the table.
- Slide the 3 colored papers out and lay them in a line on the table with the colors facing up.
- Ask a friend to think of one of the colors. Really concentrate on it.
- Say Kalamazoo and then ask him to tell you the color he chose.
- If he says "BLUE", put the papers back into the envelope, fold the flap over to close it and then hold up the front and have your friend read what is written there.
- If he says "RED", turn over the Green, then Blue, then Red papers and have him read what is written.
- If he says "GREEN", pick up the envelop, and spread it open so he can see the paper inside. Ask him to take it and read what is on it as you pick up the papers and put them away.
- Immediately retrieve the materials and put them away before they can be inspected. Do not do the same trick twice.
Find the Card
deck of cards
- Separate the deck so all reds are together and all blacks are together. Place the reds on top of the blacks.
- Fan out the deck face down on the table making just the bottom half of it spread out and the top half still piled. This will encourage the person to take a black card.
- Tell him to look at his card and show it to his buddies.
- While he is checking out his card, stack the deck back together.
- When he is ready, fan the deck again, this time trying to keep the bottom half stacked and fanning the top half.
- Have him slip his card into the deck wherever he wants.
- Stack the deck and ask him to cut the deck once.
- Stack the deck after the cut and ask someone else to cut it also.
- You can have it cut quite a few times without much chance of it going wrong.
- When ready, pick up the deck and look through it for a single card that is a different color from those around it. This is the card.
The Best Card Trick Ever
This is the card trick I learned in 6th grade and have used it ever since. It's easy but very entertaining for those that have never seen it before. And, it takes awhile to figure out.
deck of cards
- Count 21 cards off the top of the deck and give them to your friend.
- Ask him to shuffle and cut them as much as he wants. Tell him this is his only chance to touch the cards so he needs to do a good job.
- Have him give you the stack of cards face-down when he is finished.
- Deal the cards face-up into 3 columns of 7 cards each with each card in a column overlapping the one under it so they are easy to scoop into a stack keeping the cards in order. Put the rest of the deck aside.
- Point to each column telling everyone this is column 1, column 2, and column 3 from left to right.
- Ask your friend to select in his mind one card. Have him tell a friend if he wants.
- Have him tell you if the card is in column 1, 2, or 3.
- Whichever column he chooses, stack up each column and put his chosen column between the other two to create one large stack. At this point, you know his card is between the 7th and 15th cards.
- Now, lay the cards out in 3 columns again just like the start. The important thing to remember is to lay the cards out by row from left to right! This arranges them correctly.
- Ask which column the card is in now and repeat the stacking and column creation.
- Ask which column the card is in. At this point, you know the card is the middle card in the chosen column! Believe me. :-)
- When you recreate the stack putting his column in the center, his card is now the 11th card. The next part is the fun part of the trick.
- Deal out the cards into 4 card stars facedown. You will make 5 stars with 4 cards and the last one will have a 5th card in the center on top. As you deal out the cards, count them in your head so when you lay down card number 11, you know right where his card lays.
- Ask him to point (do NOT touch or the magic will evaporate) to 2 piles of cards. If he points to the pile his card is in, then swipe away the other piles. Otherwise, swipe away the piles he pointed to.
- Repeat asking him to point to 1 pile. Take away the pile or piles that do not contain his card.
- You may have to do it a 3rd time until there is just one pile left.
- Spread the cards apart a bit and have him point to 2 cards. Remove either the ones he pointed to or the others, leaving his chosen card behind.
- Repeat until only his card is left.
- Just sit there and wait for him to ask, or take all the other cards and shuffle them back into the deck and get up and go get a drink, or think of some other dramatic way to flip the card over.
How Many Cards Moved?
deck of cards
- Retrieve these cards of any mixture of suits - A, 2, 3, 4, 5, 6, 7, 8, 9, 10, Jack or Joker
- Arrange the 11 cards on the top of the deck as 6, 5, 4, 3, 2, A, J, 10, 9, 8, 7 with the 6 on the very top.
- Deal the top 11 cards from the prepared deck face-down into a row on the table going from left to right so the 6 is on the far left and the 7 is on the far right.
- Explain that you want your friend to move as many cards as he wants from his left to right, one at a time, while you turn your back.
- After he has moved the cards, turn back around and turn over the 7th card from your left. That number is the number of cards he moved! If it is the Jack, he tried to fool you by not moving any cards or by moving all of them.
deck of cards
- Separate a deck into two piles, those whose printed numeral has a ROUND top ( 2 6 8 9 10 Q ) and those whose printed numeral has a FLAT top ( 3 4 5 7 J K A )
- Put these two piles back together into one deck. It looks mixed up, but isn't really. It helps if you have the Ace of Spades as the bottom card for the Flat top group so it is easy to find in the middle of the deck.
- Divide the deck in half, giving each half to a friend. (this is where having the Ace of Spades as the center card helps.)
- Turn your back and ask them to each shuffle and cut their own decks as much as they like and let you know when they are finished.
- When they are ready, tell them to fan out their decks and pull one card from the other person's deck.
- Tell them to look at the card and then insert it anywhere into their own deck.
- If they would like, shuffle and cut their own decks as much as they'd like and tell you when they are done.
- When they are ready, turn around and have them give you their decks. Put the two decks together into one.
- Fan out the cards in your hand looking for the ROUND top number in among the flat tops and the FLAT top number among the round tops.
- Pull out the two cards and toss them on the table.
deck of cards
- Pull out all the Kings, Queens, and Jacks from the deck.
- Arrange them into 2 piles - the KH, JC, KS, QD, QC, JD in one pile, the QH, KC, JH, QS, KD in the other, The JS is not needed.
- Put the 5-card deck on your lap under the table and the 6-card deck on the table.
- Deal the 6 cards face-up into a row on the table.
- Have someone choose in their mind one of the cards and have him tell a friend.
- Ask them to really concentrate on that card as you collect all the cards into a deck and bring it under the table.
- While the cards are under the table, concentrate very hard, making mental contact with the two friends.
- Switch decks with the 5-card deck and deal out all the cards, stating that their card has disappeared!
- (of course ALL the cards have disappeared so it doesn't matter which card they chose. Don't do it twice!)
deck of cards
- Pull out all the Aces into a pile, all the Kings into another, and the Queens, and Jacks. Tell the audience that these are the King, Queen, Jack, and Ace patrols getting ready for a campout.
- Flip the piles over and stack them one on top of the other into one deck.
- Say that the patrols decided to mix it up and have one member of each patrol in a tent. Lay down (face-up) each Ace in a row. Then, on top of the Aces, flip the Kings, then the Queens, then the Jacks. Now, there are 4 piles with one of each card.
- Flip the piles over and say the Scoutmaster called lights out and its time to sleep.
- Pick up each pile, making a stack of the cards.
- Cut the cards 3 or 4 times saying that during the night the scouts tossed and turned and had a terrible time sleeping.
- Deal the cards one at a time into 4 piles (face-down) and say, "But morning finally came ..."
- Flip the piles of cards over as you say, "and they all woke to find themselves back with their own patrol!"
Finding the Aces
deck of cards
- Pull out all the Aces into a pile face-up.
- Shuffle the remaining cards and deal the top 3 cards into a row face-down.
- Deal out the remainder of the cards face-down onto these 3 cards until there are 14 on each. Then, put one more card each on the middle and right piles and put the remaining 4 cards face-down to the right as a 4th pile.
- You have 4 piles and the 4 Aces.
- Have a friend take the first pile and shuffle it all he wants, then put it back face-down.
- Have him place the top Ace face-down on top of the first pile.
- Have the friend shuffle the 2nd pile all he wants and then put it back.
- Have him take any number of cards from the 2nd pile and place them on top of the 1st pile.
- Have him place the 2nd Ace face-down on the 2nd pile.
- Have him shuffle the third pile, replace it, and take any number of cards and put them on the 2nd pile.
- Have him place the 3rd Ace face-down on the 3rd pile.
- Have him shuffle the 4th pile and place it face-down on top of the 3rd pile.
- When he is finished, place the 3rd pile on the 2nd, then the 2nd on the 1st. You still have the 4th Ace sitting aside.
- Now announce that even though the Aces are well hidden throughout the deck, you will bring them to the top.
- Deal 2 cards face-down next to each other. Continue dealing face-down alternating between the two piles until all cards are dealt.
- Pick up the left-hand deck and deal the top card from it onto the remaining deck. Deal the next card down where the left-hand deck had been. Continue dealing out cards alternating between the two piles.
- Repeat picking up the left-hand deck and dealing out the cards, first card on top of the remaining deck, until there are only 3 cards left.
- Dramatically show that these are the 3 missing Aces!
A Number from 1 to 9
deck of cards
- Give the deck of cards to a friend and turn your back to him.
- Let him shuffle the deck as much as he'd like.
- Tell him to think of a number from 1 to 9 and then deal out that many cards face-down into a pile quietly so you can not hear them.
- Tell him to now search through the rest of the deck for a card that matches the number he chose. So, if he chose 5, find the 5C or 5H or 5D or 5S.
- Once he finds his card, tell him to put the deck face-down on the table and put his card face down on top of it.
- Then, place his dealt out pile on top of the deck.
- Now, have him deal off each card from the top of the deck face-up on the table and announce the name of each card as it is dealt.
- Starting at zero, mentally count the cards he deals. When he calls a card that matches the number in your head, that is his card, so memorize it. If he calls 5 Spades when you are at '5', then that is his card.
- Let him continue to deal out cards up to 10 and then ask him to stop.
- Tell him his card was the 5 of Spades, or whichever matched.
- Then, turn around for the applause. :-)
More Cub Scout Information to Use:
Cub Scouts - information about the cub scout program
Bobcat Info - any grade
Tiger Info - 1st grade
Wolf Info - 2nd grade
Bear Info - 3rd grade
Webelos - 4th and 5th grade
To help make your Cub Scouts of America program the best it can be, use these pages to find good stuff:
Monthly Cub Scout Themes
Activities - find age-appropriate, advancement-supporting activites for Cub Scouts
Games - den or pack games just right for 1st through 3rd graders
Projects - fun projects for cub scouts to use to do a good turn
Recipes - easy recipes
Cub Scout Skits - skits that Cub Scouts love
Songs - songs to liven up a cub scout den or pack meeting
Stories - choose stories that cub scouts will enjoy and understand
Awards - see what awards are available to all cub scouts
Cub Scout Academics & Sports - extra recognition opportunities
Pinewood Derby - advice for the big race
Yukon Jack on Starting post
Yvonne N P on Outdoor Activity award
New Advanc Chair on Merit Badges
Mark Kirkendall on Project Search
Tamara on Project Search
Kim on Uniforms
Dwight S on Outdoor Activity award
Dawn on Webelos Super NOVA award
Anna R. on Arrow of Light
Contest - Ask a Question - Add Content - scout software
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I mentioned Wolbachia in my parthenogenesis post; here I will talk about it in more detail, because it’s a really cool parasite, as well as in the center of much research nowadays.
Wolbachia are Gram-negative ⍺-Proteobacteria (order Rickettsiales, family Anaplasmataceae; Dumler et al.,2001), systematically split into 6 supergroups, lettered from A to F (Lo et al., 2002). Those in supergroups C and D have coevolved with their nematodan hosts and are mutualistic (Bandi et al., 1998). Those in supergroups A and B are parasites in arthropods. Morphologically, they come in two groups: ~1µm large rods, and ~0.5µm diameter coccoids. They can be found in a membrane-bound vacuole in the host cell’s cytoplasm (Yen & Barr, 1971); in arthropods, this will be in egg and ovarian cells (see diagram above, from Weiss et al. (2009)), although it has also been found in other tissues, including the Malpighian tubules, and muscle and nervous tissue (Stouthamer et al., 1999). When found in sperm cell cytoplasm, it disappears during spermatogenesis (Clark et al., 2002).
Note that in this post and in the literature, what’s referred to as Wolbachia is Wolbachia pipientis. Before DNA sequencing, when bacteria were classified morphologically, W. persica and W. melaphagi were also thought to be related, but these have nothing to do with the famous Wolbachia (W. persica is a γ-proteobacterium, W. melaphagi is a rhizobacterium).
Being present endosymbiotically in 20+% of all known insect species (Werren & Windsor, 2000), and given that some estimate that over half of all insect species are infected by it (Hilgenboecker et al., 2008), it’s obviously extremely successful, although it should be noted that the virulence of each Wolbachia strain is different (Min & Benzer, 1997), and each species is has a different susceptibility. In Australian and Panamanian fig wasps, Wolbachia has been found to infect over 70% of all species (Haine & Cook, 2005). A host database can be found here.
The biggest reason for its success is its parasitic effect: it’s always transferred maternally in the cytoplasm of the eggs, and in order to ensure maximal transfer rates, it alters its host’s reproductive abilities to favour their own reproductive success (Charlat et al., 2003). It has four ways to do this, in contrast to the one or two ways present in other such parasites:
- Inducing thelytokous parthenogenesis in haplodiploid organisms, e.g. hymenopterans, thus making sure that only female offspring are produced (e.g. Stouthamer et al., 1993). It occurs by forcing the fusion of the two nuclei of the first mitotic division (Huigens et al., 2000).
- Feminizing males, i.e. males reproduce as females (Wilkinson, 1998). The mechanism for this varies; for example, in the pillbug, Wolbachia blocks the formation of the androgenic gland which produces the masculinising androgenic hormone (Martin et al., 1990).
- Killing males, either in the embryonic stages or later (Hurst, 1991). The mechanisms of action are still largely unknown, although recent pioneering research by Riparbelli et al. (2012) shows that it happens due to Wolbachia (purposely?) messing around with male chromosomes at various stages in development, leading to defective embryos and death.
- Inducing cytoplasmic incompatibility (CI) between infected males and uninfected females (Hurst & Werren, 2001), causing sterility (Bordenstein & Werren, 2007) and eventually reproductive isolation within a species.
The latter was the first of Wolbachia‘s effects to be discovered, by Yen & Barr (1971) in mosquitoes; Wolbachia itself was first spotted by Hertig & Wolbach (1924). What happens is that the sperm enters the cell as normal, but its chromosomes don’t decondense and fuse with maternal chromosomes due to a delay in the breakdown of the nuclear envelope in the male’s pronucleus (Tram & Sullivan (2002); see Landmann et al. (2009) for more molecular details), and so can’t enter the first mitosis, meaning they get discarded (Lassy & Karr, 1996). The embryo either becomes a haploid female (in haplodiploid organisms), or it dies. In evolutionary theory terms, in a population susceptible to Wolbachia-induced CI, uninfected females become more unfit, therefore giving a fitness advantage to infected females (Bourtzis et al., 2003), hence explaining how the phenotype persists.
Which manipulation happens cannot be predicted even within the same species, as demonstrated by Hornett et al. (2008), who found that in their North American Hypolimnas bolina butterflies, Wolbachia was a male killer while in Southeast Asian ones, it was a CI-inducer. The difference came from a dominant allele in the SE Asian butterfly genome suppressing the male-killing. Charlat et al. (2007) showed that such a mutation can become fixed in a population in under 10 generations due to the extreme selection pressure to maintain a decent sex ratio, a testament to the ecological power of Wolbachia.
However, in some cases, Wolbachia can be so prominent that the entire affected population gets a very female-biased sex ratio, which in the case of the butterfly Acarea encedon has led to females reversing their sexual roles and behaving like males (Jiggins et al., 2000). In other cases, Wolbachia is counteracted by other elements, for example the B chromosome in the parasitoid wasp Trichogramma kaykai which turns Wolbachia-feminised males back into regular males (Stouthamer et al., 2001).
It must be mentioned that part of their success is their ability to be transmitted horizontally across different species (Raychoudhury et al., 2009) – they aren’t host-specific. This has been shown as happening through parasitoids (Heath et al., 1999) or through the environment (e.g. by sharing a common food source (Huigens et al., 2000)). This also means that a single individual may have multiple Wolbachia species (or, better said, strains) co-existing and mingling inside it; the largest number I know of is eight, in the fire ant Solenopsis daguerrei (Dedeine et al., 2005). Note that this ability only comes in the arthropod-associated Wolbachia, whose genomes are more plastic, with recombination and phage-derived elements (Wu et al., 2004), none of which are characteristics present in their nematodan counterparts (Foster et al., 2005).
One study also reported the possibility of Wolbachia having transferred part or all of its genome to its hosts, albeit with only 2% of the genes able to be transcribed and none of them having any described effect (Hotopp et al., 2007).
The negative effects of Wolbachia are obviously of great interest for biocontrol of pests and disease vectors. For example, Alam et al. (2011) discuss the possibility of using Wolbachia to control the tsetse fly, a vector of trypanosomiasis; Atyame et al. (2011) do the same for mosquitoes. Such biocontrol would work by allowing a chosen genotype to dominate the population by infecting the undesired genotype (e.g. Xi et al. (2005)), or by shortening lifespans to prevent sexual maturity (e.g. Moreira et al. (2009)). Wolbachia can also lead to population bottlenecks with very few individuals becoming able to reproduce (Nice et al., 2009), which is another way to control a pest population.
In non-pest studies, Wolbachia leads to increased susceptibility to parasitoids in Drosophila (Fytrou et al., 2006). It also leads to a less effective immune system in the pillbug Armadillidum vulgare, as seen by a lower density of haemocytes and higher density of bacteria (Bracquart-Vanier et al., 2008). These would be other avenues for pest control if a similar effect is seen in pest groups.
Positive Effects of Wolbachia
Interestingly, the effects of Wolbachia aren’t all negative. In the Cimicidae (bed bugs), Wolbachia is a mutualist; getting rid of it with antibiotics reduces the amount of food the host gets (Hosokawa et al., 2010). In mosquitoes, Wolbachia was found to boost their immune system and cause resistance to dengue virus (Bian et al., 2010). Pinto et al. (2012) describe how this happens at the genetic level. This is another potential use of Wolbachia as a biocontrol agent for disease vectors. In Drosophila, Wolbachia has been shown to confer resistance to several RNA viruses (Teixeira et al., 2008). In a Drosophila lab culture, Weeks et al. (2007) showed that the Wolbachia went from being a parasite to being a mutualist within two decades.
At the extreme end, nematode-infecting Wolbachia are needed for nematode development and fertility (Foster et al., 2005), so Wolbachia antibiotics could be used to control their populations (Taylor & Hoerauf, 1999). This is useful knowledge, given that some of the affected nematodes are vectors for very serious diseases like elephantiasis and onchocerciasis. It’s a similar story with the wasp Asobara tabida, wherein no ovocyte can even be produced when Wolbachia isn’t present (Dedeine et al., 2001) because its absence promotes excessive apoptosis in the ovarioles (Pannebakker et al., 2007).
Other effects of Wolbachia on sexual physiology have been documented, for example an increase in sperm competition in Tribolium beetles (Wade & Chang, 1995), or changes to the spermathecal duct in female Allonemobius crickets (Marshall, 2007).
An interesting point can be made about the process of molecular evolution in Wolbachia. As I said in the introductory paragraph, Wolbachia is a nematode mutualist and arthropod parasite (generally speaking). One of the Wolbachia genes involved in interaction with the host is wsp, which codes for cell membrane proteins. It was found to be undergoing divergent selection when it is in a parasitic relationship, but not when it’s in a mutualistic relationship (Jiggins et al., 2002). This is in line with what we expect: wsp being involved in host recognition means it theoretically should experience heightened evolutionary rates, and this is confirmed by the empirical data.
Many, if not all, negative and positive effects of Wolbachia have evolved by natural selection in order to maximise the transmission of the strain, either by allowing the bacterium to survive in the host (depressed immune system), or reducing competition by blocking the transmission of other pathogens (as Teixeira et al. (2008) suggest for the viral resistance effect). By extension, this means that parthenogenic arthropods aren’t expected to be Wolbachia hosts, since the manipulations are useless there. In terms of evolutionary theory, they can be treated as nothing more than selfish genetic elements.
When I first heard of Wolbachia, my intuition was that it played a sizeable role in speciation, since it promotes reproductive isolation, or by selecting for subdivided populations (Hatcher et al., 2000). Some analyses showed it not to be true (Rousset & Raymond, 1991), but more and more recent studies are supportive of the idea (Bordenstein, 2003), so it’s accepted as a cause of speciation. It definitely has been demonstrated (Thompson, 1987), and in some cases has also induced rapid speciation (Bordenstein et al., 2001).
On a general evolutionary synthesis level, Wolbachia is pretty interesting as a very recognisable case of inheritable symbiosis, one of the few proper examples that lend credence to the view that symbioses are a driving force behind evolution.
Milder effects of large-scale Wolbachia infection and sex ratio-skewing include altering dispersal ability – many insects have dispersing females and non-dispersing males, or vice versa. On a more influential level, there is also evidence that they can play a role in sexual selection (Jiggins et al., 2000), since sexual conflict gets reduced when levels of polyandry fall (Arnqvist & Rowe, 2005).
Wolbachia can sometimes present a methodological stumbling block for molecular phylogenies based on mitochondrial DNA, since mtDNA will also hitchike maternally, favouring the maternal mtDNA haplotype, eventually leading to the entire dataset being worthless; see Ballard & Rand (2005) for more information. However, Arthofer et al. (2010) tested this idea using infected bark beetles and found no significant effect from Wolbachia, so it is still unsure just how significant this slight inaccuracy is.
Where it is definitely a problem is in barcoding initiatives using mtDNA. For successful barcoding of a species, a stable molecular marker needs to be used that is guaranteed not to vary across individuals, populations, or ecomorphs. However, there are some studies that show that Wolbachia causes divergences in mtDNA sequences even among individuals of the same species, e.g. in the butterfly genus Hypolimnas (Galtier et al., 2009). The reason for this is that mitochondrial genes are transferred only maternally, so only the maternal set plays a role in evolution. Given the ubiquity of Wolbachia, this is definitely a large problem that should be studied carefully before proceeding with mtDNA barcoding.
Wolbachia alone can’t be cultivated, but it is possible to keep a Wolbachia line using host cell lines (Noda et al., 2002), so experimental evolution studies are possible with them.
For the entomologists among you, make sure to check any colonies for Wolbachia infections, as they could invalidate your results, especially if you’re doing population biology. They can be gotten rid of using any antibiotic. I hear that tetracycline is recommended; if that’s not possible, high heat is enough, since Wolbachia is sensitive to temperature. If you’re sequencing your insects as well, using DNA from the legs is probably the safest way to avoid getting contaminating Wolbachia DNA amplified (this is standard procedure anyway).
Other symbionts that alter the reproduction of their arthropod hosts include Buchnera and Cardinium – but I’ll leave them for other posts.
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Research Blogging necessities :)
R. Stouthamer, J. A. J. Breeuwer, & G. D. D. Hurst (1999). WOLBACHIA PIPIENTIS: Microbial Manipulator of Arthropod Reproduction Annual Review of Microbiology DOI: 10.1146/annurev.micro.53.1.71 |
51°15' / 22°34'
Translation of "Lublin" chapter from Pinkas Hakehillot Polin
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Encyclopedia of Jewish Communities, Poland, Volume VII, pages 13 - 38, published by Yad Vashem, Jerusalem
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Lublin (hereafter L), the largest town in Eastern Poland, is the main town of the region and province. It stands on the banks of the River Bystrzyca. In the 12th century it was a fortified settlement on the trade route leading to Ruthenia. In the 13th century it was destroyed several times during invasions by the Tatars, Ruthenians and Jadvingians (a Lithuanian tribe). For some time the settlement was under the ægis of the Princes of Halicz. In 1317 L was granted the Magdeburgian privileges of a town by King Wladislaw Lokietek. His son, Casimir the Great, erected a wall around the town. It was granted the privilege in 1385 of free trade with the Lithuanian Principality and in 1392 the right to store in the town goods in transit - and these advantages contributed to its rapid development. In 1474 L was the seat of the king's representative (the Woiwoda). It reached the zenith of its prosperity in the 16th and beginning of the 17th centuries, when its markets attracted merchants from the whole of Poland and even beyond. The same period also saw the expansion of various trades.
In the 16th century L was also a political centre, where the nobility held its meetings; in 1578 it became the seat of the Crown Court and the Appeal Court of "Lesser Poland", and in 1569 there convened in L the Sejm that united Poland and Lithuania under one crown (this was the celebrated Polish-Lithuanian "Unia"). Resident in the town were writers and poets renowned at the time (Kochanowski, Rej and others), and a printing works and paper factory were established.
In the 16th and 17th century the population of L was more than 10,000, and embraced some suburbs. In the period of the "deluge" (the wars in the middle of the 17th century waged in Poland by the Cossacks, the Swedes and the Muscovites) L suffered heavy damage to life and property. Repairs were not undertaken for some time, and only in 1780 did work on reconstruction begin. In 1795 L passed under the sovereignty of Austria, in 1809 it formed part of the "Principality of Warsaw", and from 1815 until the rise of independent Poland in 1918 it was included in Congress Poland.
In the course of the 19th century L developed into the metropolis of the eastern part of the Polish Kingdom. Its population increased from 7,000 at the beginning of the century (during the Napoleonic Wars) to 56,000 at its conclusion. During the century it was under the "Jurisdiction of Wieniawa" (a suburban noble estate with autonomous legal status). Only late in the century did L start to develop a modern economy. After being linked to the rail network in 1877 existing industries developed and new ones were established. Concomitant with this there was cultural development and cultural and educational institutions were founded. L contained some well-known Polish writers, such as Kraszewski, Vincenty Paul and Boleslaw Prus, and many figures of the Polish national underground movement chose to settle there.
The 20th century saw no radical changes in L's economic structure, and the town remained the centre of an agricultural area. World War I had a considerable stagnating effect on local industry, and its recovery between the two world wars was protracted. In the 20s and 30s new plants connected with the aviation industry were established, as were factories for pipes, screws, tobacco and chicory (coffee substitute). The town's public services (sewage, waterworks, electric power, slaughterhouse and the like) were also developed. Nevertheless, during the whole inter-war period L was in economic decline, or at best stagnant. The number of unemployed went on increasing and with it the number of demonstrations and strikes. The process of impoverishment affected broad spectra of the population. The situation was particularly bad during the crisis of 1929-33, and even in 1933 the number of unemployed (12,000) surpassed the number in employment. There was virtually no natural increase of population in the years preceding the Second World War.
During the 19th and 20th centuries L was the scene of activity by political movements and youth organisations, both open and clandestine. In 1918, after Poland gained its independence, a council of Polish representatives assumed political control of the town and its environs.
On the very first day of World War II L suffered heavy bombardment. During the German occupation the leading Polish intellectuals were liquidated. Many people were imprisoned in the town's fortress on political grounds. In 1941 the first batch of 160 prisoners was despatched to the concentration camp at Ravensbrück. In 1941 too, an extermination camp was established at Majdanek, near L. It served at first as a camp for 2,000 Soviet prisoners of war. In Majdanek some 360,000 beings met their death, most of them Jews from Poland, but also from other parts of occupied Europe.
During the German occupation whole quarters of the town were destroyed, as was most of the infrastructure (such as roads and bridges). L became an important centre of the Polish Resistance.
On July 24th, 1944, L was liberated by the Red Army, and for seven months was the capital of the new Poland. The process of recovery was rapid: ruined factories were restored and new ones built, and L again became an important cultural centre. New institutions of higher education were set up, including the University of Marie Sklodowska Curie, as were various categories of secondary schools. The Theatre and the Philharmonic Orchestra resumed their activities. The Lupaczinski Library was expanded to contain 200,000 volumes. However, in the 60s stagnation and even recession set in, as was the case in the rest of "Socialist Poland". The crisis reached its peak in the late 80s, and even continued into the beginning of the new Polish Republic that came into being in 1990-91, following the collapse of the "Polish People's Republic".
The first Jews in L are mentioned in documents from the 15th century, though there may possibly have been some there earlier. A few Jews lived in the town itself, with the consent of the authorities, despite the royal edict of 1535 that forbad Jews to live in L. In the 15th and 16th centuries many Jews settled in the town and established a Jewish quarter to the north and north-east of the palace and the fortress. For the privilege of living there the Jews were obliged to pay a tax to the fortress administrators. In 1555 the boundaries of the Jewish quarter were extended; the Governor of Lublin, Stanislaw Winczinski, with the approval of the King of Poland, Sigismund August, granted the Jewish community three large plots of land, One of these was to contain the slaughterhouse and butcher shops; the second, on the palace hill near the old cemetery, was to serve as an additional burial-ground; and the third, on the slopes of the palace hill, was set aside for religious purposes. In return for these plots the community agreed to pay an annual tax of 12 marks and an annual sum for candle grease.
In 1557 Governor Winczinski presented Dr. Yitzhak Maj, a well-known doctor in L, a large plot (which also had a pool) and permitted him to build on it. In 1566 a synagogue was erected there in the name of the Maharashal (Rabbi Shlomo Luria), and the following year other buildings were added, in one of which a yeshiva was opened. In 1603 the merchant Yaakov ben Aharon received from Governor Firlay a plot on which to erect buildings and shops. The Jewish quarter continued to grow, from 24 houses in 1550 to 100 at the end of the century. The main street was Sieroca Street, where the activities of the community were concentrated. In addition to the Maharashal Synagogue and the Bet Midrash, there was another prayer-house, in the name of Avraham ben Chaim, one of the leaders of the Council of the Four Lands*, and the study house of the "Seer of Lublin". A third synagogue was built in Podzamcze Street - the "Leifer", associated with the name of Shaul Wahl, of whom it was said that he ruled Poland for one day, but gave up power in favour of a Pole. In 1638 Tsvi Hirsch, an agent of the royal court and a leader of the Council of the Four Lands, built a synagogue in the name of his father, Moshe Doctor ("Doctor-Shul"), this too in Sieroca Street.
In time the Jewish quarter became too restricted and some of its inhabitants tried to circumvent the rule of "de non tolerandis judais" (not to tolerate Jews within the confines of the city) implemented by watchful citizens, with a view to moving to other parts of the town. Few, however, were permitted to do so, and these only for a limited period. In 1655 a fire broke out in the Jewish quarter and most of the houses went up in flames. Reconstruction took many years. As mentioned above, after the "deluge" (the wars of 1648-60) L lost its privileged status. The Council of the Four Lands, which until then usually met in L, moved in 1680 to Jaroslaw, which was developing at the time and was an important centre of trade and markets.
In the 18th and the second half of the 19th century also, the Jews continued their attempts to be allowed to live outside the Jewish quarter, but generally without success. In 1761 Jews who were there without authorisation were expelled from the centre of town and, with a few fortunate exceptions, were compelled to return to the crowded conditions of the quarter. This state of affairs continued until the Tsar's edict of 1862 annulling the restrictions on Jewish dwellings in the Kingdom of Congress Poland.
Thanks to its status as an important commercial centre as early as the 15th century and its position at the crossroads of the trade route from Western Poland to the North and South-East of the country, the markets of L attracted large numbers of Jewish merchants from all corners of Poland and beyond. The Jews traded in almost all the wares bought and sold on these occasions. In 1530/31 Jewish merchants transported through the auspices of the Chamber of Commerce in L 136 barrels of wine, 503 oxen, 15,000 ells of cloth, 47 rolls of silk, 38,297 furs and hides, 5 wagons of iron, 82 stone of steel, 22,170 ells of fabrics, and other wares. In 1584 Jewish merchants from Krakow brought to the market in L 185 wagons loaded with goods. The communities of Poznan and Krakow even had in L brokers, whose task was to settle differences between the merchants during the markets. For their part, the merchants of L maintained contact with important merchants outside Poland.
At the end of the 15th century and the beginning of the 16th the name of the merchant Yoska Szachnowicz of L (the father of Rabbi Shalom Shachna Hacohen) was well-known for his trading and leasing activities. Yoska was the King's excise officer and in charge of customs in Reisen and Podolia.
The markets of L, where Jewish notables from all over Poland gathered, became in time the core around which the Council of the Four Lands grew. However, before they achieved commercial status the Jews of L had to overcome many difficulties. In 1518 King Sigismund I acquiesced in the request of the citizens and town council, and ordered his representative in L, Jan of Filicia, to curtail the trading rights of the local Jews. In 1521 the king issued an edict forbidding Jewish pedlars to purchase produce brought into town by the local farmers. Two years later another order was published with a view to placing the trading rights of the L Jews on an equal footing with those of Lwow, Poznan and Krakow. This order contained a set of restrictions, mainly a ban on retail trading by Jews and a curb on their wholesale activities. The Jews of L were also forbidden to have shops, and were only allowed to set up temporary booths in their own quarter. Christians were forbidden to lease shops or stalls to Jews, infringements being punishable by fines. In spite of this, the Jews managed to rent stalls and shops during the annual markets.
In 1531 there was an action brought by the Municipality of L against the Jews concerning payment of tax on imports of woven goods. The king ruled that the Jews had to pay this tax, despite the fact that non-Jewish citizens were exempt. Contrariwise, there were instances where the Jews were freed of various taxes by royal order; for example, in 1530 the Jews were exempted from paying road tax.
Generally speaking, the Jews strove to arrive at a compromise with their non-Jewish townsmen. The Christians in the Podzamsze quarter held a monopoly on the right to produce and market wines, beer and mead; and in 1590 the king forbad the Jews to engage in these activities. They, however, managed to circumvent this ban through an arrangement whereby, cash in advance, they could sell spirits purchased from the Christians. In return for this favour, the Jews agreed to make an annual gift of 35 ells of finest quality cloth to each member of the town council, and twice a year a pound of pepper to the town clerk and the seven members of the council. The Jews were forbidden to sell bread and flour to non-Jews. This agreement was renewed by the two parties every few years, with the approval of the king. From time to time new edicts and restrictions were imposed on the Jews. In 1650, as a result of Jesuit pressure, Jewish apothecaries were forbidden to prepare or sell medicine. And four years later, in 1654, Jewish musicians were banned from appearing in public without special permission from the town authorities.
In their battle against Jewish trade the other citizens employed not only legal means, but also violence. Now and then they would burst into the Jewish quarter and kill, and loot property and wares. Such an incident occurred in 1623. There was also violence for its own sake. In 1646 pupils of the ecclesiastical "colleges" (also known as "Dzakim") fell upon the Jewish quarter. In the course of this disturbance eight Jews were killed, many injured, and 20 houses emptied of their contents. In 1598 a rumour of blood libel spread through the town; three Jews of L were accused of murdering a Christian child for ritual purposes, and the mob tried them and slaughtered them brutally. Members of the Jesuit Order in L set up a printing press and distributed leaflets with virulent accusations against the Jews.
Jewish artisans too were forced to wage a hard battle against their non-Jewish competitors for the right to engage in their trade. Even before the "deluge" the Christian artisan guilds enjoyed privileges that entailed bans or restrictions on Jewish business. In 1535, for example, the Christian tailors' guild obtained from the king an order banning Jewish tailors from engaging in their trade, and in 1595 the bakers' guild obtained a similar decree from King Sigismund III forbidding Jews to bake bread or trade in flour. This devout king also granted privileges to the guilds of carpenters, glaziers and ropemakers, according to which the guilds were excluded from accepting as members unbelievers, heretics and Jews.
The municipality, the official patron of the arts, also approved many measures directed against Jewish artisans. In the 17th and 18th centuries Christian artisans, with the approval of the municipality, adopted regulations designed to forbid the Jews to engage in various trades, or to restrict them to marginal activities. The Christian artisans also obtained from the authorities a ban on Jewish merchants supplying various trades with raw materials.
However, in not a few cases the Christian organisations waived some of the prohibitions and restrictions in return for money; a typical example of this was the agreement of 1689 with the guild of weavers of gold and silver threads, by which three Jewish weavers were permitted to work in eight workshops, and 13 Jewish brokers and suppliers in the branch were allowed to buy and sell raw materials - apart from Christian ritual articles - exclusively to non-Jewish artisans. In return the Jews consented to pay to the guild 100 guilders four times a year and to accept supervision of their workshops by the guild. Clearly the Jewish artisans could not put up with these numerous restrictions and endeavoured to circumvent them by all means possible, both legal and illegal. In 1600 a group of Christian tailors seceded from their guild and started a separate guild in Podzamcze (a residential area open to Jews). The new guild accepted 24 Jewish tailors on payment of an annual fee. Since, however, the number of active Jewish tailors far exceeded those accepted by the guild, and since this majority also had to earn a living, many of them plied their trade illegally. This led to friction and violence. In 1781 a new agreement was signed by which the number of Jewish tailors permitted to work was increased to 42; but this new arrangement also included the old restrictions and obligations incumbent upon the Jews. In 1792 the parties concerned drew up yet another agreement, but neither did this put an end to dissension. The Christian guild continued to complain that the Jews were in arrear of payments and that they employed Jewish tailors from other towns. In 1805 the two tailoring guilds of the town amalgamated, and 15 years later the number of Jewish members had surpassed that of the Christian. One restriction though still remained: Jews were not eligible for election to the guild council. But in 1832 this disability was also removed, and two Jews were elected to this body.
Similar struggles were waged by other Jewish artisans. For example, in 1595 the Christian bakers were granted by the king's representative the privilege of trading in flour and baking products - this was denied to the Jews. This right was renewed in 1611 by Sigismund III, and again several times by successive Polish kings. Yet after a difficult struggle by the Jewish bakers a breach was made; it transpired that the Jews were allowed to bake products for their own use, and they also exploited this paragraph with extra-legal activity.
The attitude of the royal authorities in the town to Jewish artisanship was far more obdurate than towards Jewish trade. The reason for this was mainly fiscal: Jewish trade constituted an important source of revenue for the king and the town, and the authorities were therefore prepared to compromise and allow Jewish trade, albeit with various restrictions. Despite such difficulties and curbs, however, Jewish workshops in L increased their activity and extended into more and more branches. According to the Census of 1764/5 there were 38 artisan workshops in all, nine of which were unspecified. Over 20 years later, in 1786, there were 30 tailor-shops (only nine were Christian); 30 shoemakers (masters only, excluding journeymen and apprentices); 30 glaziers (only one Christian); an unspecified number of jewellers (nearly all Jews - only two non-Jews are mentioned); and furriers, where Jews outnumbered non-Jews by four to one. The tinsmiths and coppersmiths were all Jews, as were the bookbinders and coopers.
In addition to trade and artisanship the Jews engaged in credit transactions and leased the right to collect national and municipal monies (taxes and public revenue). Some of the tax and rent collectors and lessees prospered greatly and accumulated much capital. Joska Szachnowicz, who lived at the end of the 15th and the beginning of the 16th century, and was the main lessee of royal revenue in "Lesser Poland" and in Reisen, has already been mentioned. After his death in 1507 his wife Golda took over his business, and she was followed by her son, Pesach Joskowicz (his younger brother, Szachna, was known as the greatest of L's rabbis). Much business in credits was carried out in the 16th century by the doctor Yitzhak Maj and his brothers Moshe and Joel. The sum of these transactions amounted to 5,000 guilders, an enormous figure for those times. In 1604 Moshe Doktorowicz (son of a doctor), a leader of the community, leased tax collection of the Jews in 20 towns in the districts of Belz and Chelm.
The struggle of the Jews to earn a living was especially difficult because of the heavy tax burden. These included royal taxes, municipal taxes, and other duties and payments, both regular and extraordinary - property taxes to estate owners (juridically included in L), and the constant need to give presents to the royal functionaries, to the artisan guilds, to priests, to the Jesuit schools, and so on and so forth.
In addition to a poll-tax (one guilder for every Jew over the age of one - this was raised to three guilders in the 18th century) the Jews of L in 1640 paid to the royal treasury (via the Starustra) 700 guilders a year. To this sum was added an annual tax to the Woiwoda (300 guilders), payments for maintenance of the royal court and entourage whenever they passed through L, and a tax called "royal szos" on houses and land. The Jews of L contributed "voluntarily" to the liquidation of the kingdom's public debts - their share sometimes amounted to as much as 60,000 guilders.
The Jews of L paid a whole series of direct and indirect taxes to the municipality and to the council of Podzamce, where the Jewish quarter lay. The shopkeepers, artisans, and house owners paid duties on goods, funeral tax, tax on drinks, market tax, import tax, and so forth. In addition to all these, the Jews also paid taxes to the community - a direct tax ("sum") and indirect taxes ("maintenance", "clientele", etc.). The burden of payments on each Jewish head of family was very heavy, and many could not meet the demand.
In times of war taxes and duties on the Jews were doubled and redoubled. The Jews paid most of the enemy's impositions on a town under siege and the conqueror's ransom money, amounting to thousands and tens of thousands of guilders. In 1676 the Jews were required to pay 8,000 guilders for military expenses. In 1698 they contributed 4,000 guilders towards renovation of the town hall. Large sums were also demanded to finance the town's water supply and other services. Contemporary documents naturally make no mention of the sums paid to officials, to merchant guilds or to artisan guilds.
The municipal taxes had no foundation in law, since the Jews were not considered citizens of the town - but the municipality and the Christian organisations took no account of this. Nor did the appeals of the Jews to the royal court help. Thus, Jewish debt grew and grew, and for many years (until the beginning of the 19th century, after the division of Poland) not only was the community unable to repay the principal - it could not even pay the interest and compound interest to be paid to its creditors (the nobility, the church, and the non-Jewish citizens).
The community of L was considered the most important in Poland, by virtue of its reputation as a centre of Jewish Law and Culture. Almost all the celebrated rabbis of the time officiated for shorter or longer periods as rabbis of L. In 1475 R. Yaakov Matrident (died 1541) arrived in L, and his presence there brought added renown to the community.
The most eminent rabbis of Poland throughout time graced the rabbinical chair in L. First there was R. Yaakov Halevi (died 1541). His successor, R. Shalom Shachna Hacohen, who had lived in L since 1520, was not only L's spiritual Jewish head, but also a man of wealth and responsible to the State for the public finances of the community. Nearly all the renowned authorities on Halacha (Jewish Law) in Poland in the 16th century studied at the yeshiva he established - among them the Rem”a (R. Moshe Isserles). In 1541 he was appointed by the king Chief Rabbi of all "Lesser Poland". His gravestone in L has remained untouched. After his death in 1559 he was succeeded by his son, Israel, who set up another Yeshiva, headed by R. Shlomo ben Yechiel Luria (the Maharashal). R. Luria's books "Yam shel Shlomo" (Solomon's Sea) and "Chochmat Shlomo" (The Wisdom of Solomon) are cornerstones of Talmud and Halacha.
Rabbis Isserles and Luria were considered the leading exegists of their time. R. Luria was one of the founders of "Nusach Polin" (the Polish Version) of Judaism. After his death in 1574, the post of Chief Rabbi went to his son-in-law, the Principal of the Yeshiva, R. Shimon Wolf Auerbach, son of R. David Teyvele - he too renowned in his time. R. Auerbach officiated as Rabbi of Prague towards the end of his life, and died there in 1632. His successor in L was R. Mordechai Jaffe, author of the Levushim, a popular exposition of Jewish Law; he too was a renowned scholar and also the founder of the Council of the Four Lands. He died in 1612. At the turn of the century until his death in 1623 the L rabbinate was occupied by R. Yitzhak ben Nute Hacohen. Also Chief Rabbi of L was R. Yehoshua Falk, a pupil of the Maharashal. Another renowned rabbi was the Maharam of L, Meir ben Gedalia, who served as Principal of the Yeshiva and Rabbi of the town (died in 1616). Other Rabbis of L in the 17th century worthy of note are Shmuel Eliezer ben Yehuda Eidles, the Maharasha (1614-25); Yoel Jaffe Sirkes; Naftali ben Yitzhak Hacohen; Efraim Zalman Shor the Elder, author of "Tvuot Shor" (Shor's Harvest) (died 1634); Avraham Halevi Epstein (1638); and Yaakov ben Efraim Naftali Hirsh.
In the second half of the 17th century too there were many eminent rabbis, amongst them: Naftali ben Yitzhak Katz, son-in-law of the Maharal of Prague (died 1650); Heshel ( son of the former Rabbi Yaakov), whose name was a symbol of Jewish wisdom and perspicacity, and who was popularly known simply as Rabbi Heshel (died 1652); Aharon Shimon Szapira (died 1680); Moshe, son-in-law of the Maharasha and author of "Mahadura Batra" (The Last Edition), and his son Israel Isser (died 1692); Tsvi Hirsh ben Zecharia Mendel (died 1690); and Mordechai Ziskind, author of "Shu"t Hamaram Ziskind" (Responsa of R. Mordechai Ziskind).
The voices of the Rabbis of Lublin within Jewish Law and Tradition echoed throughout Poland and beyond. In 1587, by special permission of King Sigismund August, there were established in L a Yeshiva and a Bet Midrash. Their first principal was R. Yitzhak Maj, who was not subject to the authority of the town rabbinate.
In all matters relating to the internal affairs of the community its leaders and rabbis followed Hebrew Law. In 1631 the Woiwoida, Petro Tarlo, confirmed the right of the Jews of L to manage their affairs on the basis of their religious laws, and to try cases involving Jews only in a Jewish court. Only in cases needing clarification, or those where one of the participants was a non-Jew, was judgment made in the court of the Woiwoida (Governor) or the king.
L plays an important role in the annals of the Hebrew press in Poland. As early as 1547 sacred Hebrew books were printed there at the works founded by the Shachor family. The first printer in the town, Josef, bequeathed his printing works to his daughter Chana and her husband Yitzhak. In 1559 these two were granted the royal privilege of a monopoly on the publishing and distribution of Hebrew books. Despite this privilege, however, another press was set up near L, in Konskowola. Neither of these two printers survived long. In 1566 two Jews, Eliezer ben Yitzhak and Josef, were in turn given the royal privilege of starting a printing press, but this too had a brief life. In 1578 Klonimus Jaffe, son of R. Mordechai Jaffe, obtained permission from King Stefan Batory to set up a printing press. This works was passed on by Klonimus Jaffe to his descendants, and, with short intervals, was still operating at the beginning of the 19th century. Between 1550 and 1690 it published no less than 103 books. The presence of many printing works in L is evidence of the community's status as an important Jewish cultural centre in Poland.
L was the home of doctors renowned in Poland and beyond. At the beginning of the 16th century there was a Jewish doctor there called Yechezkiel who, in return for his services, was exempted by King Alexander Jagiello from all taxes. Another Jewish doctor, Yitzhak Maj, occupied a prominent position in the community. As mentioned above, in 1557 he was granted a plot of land on which to build a house of prayer and a yeshiva. One of his contemporaries in L was Dr. Moshe Doktor. His son, Tsvi Hirsz was a leader of the community and an agent at the courts of the kings Sigismund III and Wladislaw IV. At the end of the 16th century there lived in L the doctors Shlomo Luria (a relative of the Maharashal), author of the medical book "Luach Chaim" (Table of Life); Shmuel ben Matatyahu; and Moshe Montalto, the descendant of a well-known Marrano family from Portugal, bearer of the title "Medical Adviser to the King of France". He died in 1637, and his gravestone still stands in L. His father, Eliezer, returned to Judaism and served as doctor to the French kings Henri IV and Louis XIII. In the 17th century the doctor Chaim (Felix) Vitalis was renowned in L. He had studied In Padua and served as doctor to the court of King Mikael Wisnowiecki. This king invested him with the authority to test all the Jewish doctors. In the second half of the 17th century some of the members of the Jewish community of L went to study medicine in Padua. The best-known of these were Avraham Szapira and Moshe Izrael Polak ben Yitzhak.
Within the community there arose organisations and commitees concerned with welfare and mutual help, as well as with education and culture. In L, as in the other Jewish communities in Poland, the foremost among these was the Chevra Kadisha. This was a closed and privileged body, not open to everybody. In its hands, and through it the community, lay the authority to impose on the Jews discipline and moral precepts. It was enough for the Chevra to threaten an offender to be placed on its black list (i.e., when his time came they would not arrange his burial) for him to repent. The offences recorded in this connection were slander, infringement of the rights of possession, contact with non-Jewish legal instances, etc. Among its other tasks the Chevra took upon itself reception of the emissaries, expounders of the Torah, preachers, writers and contractual experts from many countries who came to the community in L; and it also administered the Hospice for the Needy. In the period 1690 to 1760 the records of the Chevra name guests from Italy, Turkey, Prague, Moravia (Mahren), Belgrade and Budapest, in addition to those who came from all parts of Poland. As well as the Chevra Kadisha, documents of the 18th century mention the institutions "Ner Tamid" (Eternal Light), "Bikur Cholim" (Visiting the Sick), "Rodfei Tsedaka" (Seekers of Charity), and " 7 Kruim (The 7 Guests)". Ner Tamid was to be found in every prayer-house.
The Jewish children were taught in the Cheder by private teachers. There was a Yeshiva in L, headed by the town rabbi and supported by the community council. The Yeshiva students, some of whom were not from L but who came from all parts of Poland, were supported by the "houseowners", lived in their houses, and ate daily at their tables - each day in a different house.
In the 16th and 17th centuries the Council of the Four Lands convened in L, but moved in 1680 to Jaroslaw. The Rabbi of L, Yehoshua Falk-Cohen, was the leader of the Council and the author of its set of rules. L was also the seat of the Jewish Rabbinical High Court in Poland, among whose duties was granting permission for the publication of Hebrew books. One of the latter was the edition of the Babylonian Talmud (1559-80), that was published to meet the needs of the Torah schools and other Jewish schools (cheders and yeshivas) in Poland.
In 1655 Lublin was attacked by the Cossack hetman, Zoltarnko, and the Muscovite Voivoda, Piotr Ivanovitch. The town was occupied after a long siege. The conquerors imposed a tribute of 10,000 guilders on the town, in return for which they promised not to inconvenience the citizens. They did not, however, honour this pledge. First of all, they burst into the Jewish quarter and slaughtered some 10,000 Jews, among them refugees from the surrounding area who had sought refuge in the fortified city, and plundered whatever they could find. A group of Jews wrapped themselves in their prayer-shawls and asked their tormentors to bury them alive in the Jewish cemetery (so that they would be sure to be interred in Jewish soil); the butchers acceded to this request. The names of the victims are engraved on the wall surrounding the cemetery. A description of these events is given in the book "Hesed Shmuel" (The Mercy of Samuel - Amsterdam 1699), written by Szmuel ben Dawid. The massacre took place during the Feast of Succot, and for many years it was the custom in L to mourn in this period.
The damage wrought on L by the Cossacks was not repaired for some time, and only at the beginning of the 19th century did the population return to its former level, before the "deluge", i.e. 4,000 souls. The economic derpression that followed in the wake of the Cossack aggression continued till the middle of the 19th century, when the town began to recuperate.
On the conclusion of the bloodbath, as was the custom of the times, the king promulgated a number of orders designed to alleviate the plight of the victims, and the Jews were included in these. A decree of 1669 allowed them to trade inside and outside the walls. Permission was likewise given for the Jews to acquire houses and building plots. In 1679 King Jan Sobieski replaced this decree with one that still allowed the Jews to do business within the walls, except on Christian Holydays. In 1675 the Satrusta Danilowicz issued a decree placing the Jews of L and the "little they possessed" under his protection. In 1696 the king renewed the right of the Jews to trade within the walls, like the Christian merchants.
However, these days of clemency did not last. During the "Catholic Reaction" in the first half of the 18th century, the non-Jewish citizens renewed their efforts to eject the Jews from business activity, and even to expel them from the town. There were again calls for rigorous measures against them. In 1720 King August II issued an ordinance accusing the Jews for economic and religious reasons, as it were, of mocking the Christian religion, and forbad them completely to engage in commerce; he further ordered them to leave the dwellings they had rented "by guile" from the Christians of the town.
The last king of the Kingdom of Poland, Stanislaw August Poniatowski, confirmed, however, in 1769, the "privileges" granted to the Jews of L in the past - only in 1780 to issue a decree forbidding them permanently to live or engage in trade in L and its suburbs. This situation prevailed until the town was occupied by the Austrians in 1795.
The period under review was also marked by a decline in the internal affairs of the Jewish community in L, which was enmeshed in debts and unable to pay the interest and compound interest on the loans taken, and certainly unable to repay the principal. The status of the community declined drastically, and individuals and groups began to act with an irresponsibility that bordered on anarchy. Many stopped paying their taxes to the community and avoided payment of sums the community had undertaken to pay to the authorities. Nor did they honour the code of conduct in business and crafts, with the result that unrestrained competition was rampant, and this seriously affected the livelihood of many of their fellows. These miscreants placed their trust in the hands of the estate owners near the town, on whose land they settled. The Community Council complained to King August III about this situation; and he empowered it to exercise force and even to impound possessions in order to collect the outstanding taxes. The Chevra Kadisha issued a proclamation warning the Jews of L that it would use all its power to combat disturbers of the peace, slanderers, and Jews having recourse to non-Jewish instances. Here too, as in the past, the Chevra threatened to blacklist their names: they would not be given a Jewish burial, and the protection of the nobility would not help them.
The second half of the 18th century heralded the spread of Chassidism. At its beginning the rabbis who officiated in L were "mitnagdim", like Yaakov Chaim (died in 1769), son of the great mentor R. Avraham ben Chaim, who was the leader of the Council of the Four Lands, and examined the learning abilities of the Baal Shem Tov. After R. Yaakov came R. Shaul Margules (died 1788), who was not a Chassid himself, but whose father, R. Meir Margules, belonged to the circle of the Baal Shem Tov. There is no doubt that only in the time of the latter was R. Yitzhak Horowitz, the "Seer of Lublin", able to fix his place of residence in L and to make it the "Jerusalem of Chassidism". Large numbers converged on hsi Bet-Midrash, and the whole following generation of holy men in Poland were his students. The "Seer" taught and inspired many renowned "Admorim" (Chassidic Rabbis) and Heads of Chassidic dynasties throughout Poland, among them the Heads of the Dinover dynasty, R. Tsvi Elimelech Szapira, and his son, R. David; R. Shalom Rokeach, Head of the Admorim of Belz; and many others.
Opposed to the Chassidim was the Rabbi of the Community, R. Azriel Horowicz, known as “der ayzerner kop" (died in 1819). The struggle between Mitnagdim and Chassidim weakened the community. With the death of the "Seer" the position of L as a centre of Chassidism declined.
In the period under review many different rulers held sway over the area which included L; but the legal and civil status of the Jews hardly changed. Neither the Austrians nor the authorities of the "Principality of Warsaw", which spoke in its constitution of equality for all its subjects, granted equal status to the Jews, or alleviated their restrictions. The period of a "decade", during which equal rights for the Jews were postponed "until they mended their ways", lasted till 1862. The regime of Congress Poland that followed made no changes in the legal status of the Jews in that kingdom, including the Jews of L. The so-called privileges granted by August II (in 1720) and Stanislaw August Poniatowski (in 1780), according to which the Jews were permitted to live permanently in the "Jewish town" of Podzamsze and in other quarters separated from the Christian sections of the city, remained in force until all restrictions on Jewish residence in Poland were annulled by the Tsarist Decree of June 1862.
During the Polish uprising of 1863 Jewish inhabitants of the town sided with the insurgents.
From 1765 until the 50s of the 19th century the Jewish population of L increased fivefold or more. Most of this was due to natural increase, though part was also explained by Jewish migration - legal and illegal - from nearby provincial towns. From the annulment decree of 1862 to the end of the 19th century the number of Jews in L grew from 8,747 (1857 figure) to 23,586 in 1897. An additional increase took place in the first two decades of the 20th century, numbering in 1921 37,337, or more than a third of the total population.
In the second half of the 19th century the Jews of L began to settle in the "old town" within the walls. In the course of time, as the Poles preferred to move outside the walls, the Jews formed a majority and the "old town" thus became a new ghetto.
During this period the Jews of L earned their living mainly by petty trade and artisanship. There were times when the number engaged in the latter surpassed that of those engaged in the former. The area around L was agricultural. Until 1864 most of the villagers were poor serfs tied to the soil, and their purchasing power was thus very limited. Most of them lived off their own produce and had little need to buy goods. One of the branches that depended on the peasants of the area was the production and sale of strong drinks, and here there was competition between the Jewish lessees and the declining nobility. The few implements needed by the peasants, mainly of iron, were not enough to warrant the establishment of industry or the development of trade.
The main agricultural product was corn, and trade in it gave a living to many Jews. They bought the produce at markets, or through the middlemen and distributors of the estate owners, and sold it in Warsaw, or shipped it on rafts to the port of Danzig (Gdansk), with a view to export abroad. Other important products traded in were hides, pig bristles (for brush-making), and wood felled in the surrounding forests. These goods too were sold in Warsaw or exported abroad.
The majority of Jewish merchants, however, were merely small shopkeepers or pedlars who wandered round the villages with their merchandise (haberdashery, i.e. needles, buttons, and the like, cheap cloth, etc.). In return they bought in the villages eggs, poultry, feathers, etc. The impoverished rural surroundings were the deciding factor in the miserable situation of the Jewish merchant in L and the extent of his business. No significant change took place in this situation after the emancipation of the serfs. The tiny plots of land given to the latter were insufficient to produce an agricultural surplus for sale.
Under these conditions no real industry developed in L, and the share of the Jews in what little manufacture there was remained small: of the 27 larger industrial plants only three were in Jewish hands. Most of these local industries were engaged in processing the agricultural produce.
In contrast, the number of Jewish artisans in L increased. Proof of this was the continued existence of craft branches begun in the preceding centuries.The most prominent, as far as numbers and organisation were concerned, were the tailors and furriers. Towards the end of the 19th century the salesmen were almost as numerous. There was considerable Jewish representation in all artisan branches; and in general non-Jews were in a minority in them.
In the last two decades of the 19th century three provident funds were established in L. In 1906 they embraced 1,200 members (artisans and small traders) and their equity capital amounted to 91,000 roubles. The clerks and petty traders set up a similar organisation, called "Igud Le'ezra Hadadit shel Pekidim Usocharim Zeirim Bnei Dat Moshe" ( The Mutual Aid Society of Clerks and Small Traders of the Mosaic Faith).
As stated, in the second half of the 18th century the situation of the Jewish Community Council in L deteriorated as a result of growing debt, and thus led to growing dependence on the local authorities, who were also the main creditors of the community. The interference of the Woiwoda and the Starusta in its activities was a frequent occurrence. The collection of community taxes was subject to the approval of the Woiwoda.
Within the community itself conflicts erupted that often involved non-Jewish mediation.
There was constant friction between the affluent merchants (who were dominant in the council) and the lesser merchants and artisans. Even the Rabbinate, which, as mentioned, was one of the most important in Poland, was not spared the animosity of these conflicts, and its incumbents were the subject of quarrels without end. The diagreements that arose out of the difficult financial situation grew even worse with the advent of Chassidism.
The propagator of Chassidism in L was, as stated, the Seer of Lublin, R. Yaakov Yitzhak Horowicz. During his tenure the community split into two: to begin with the Mitnagdim were in the ascendant, under the leadership of R. Azriel Horowicz, (der ayzerner kop) (died in 1819); but after his death the Chassidim grew stronger, and they determined the choice of the Chassid R. Yehoshua Heshel, Head of the Rabbinical Court at Tarnopol, as Rabbi of L. He was not, however, well received by the Jews of the town: three brothers, prominent members of the community, and Mitnagdim, instituted a campaign against him, and in 1826 he was forced to resign. From that year until his death in 1843 the Rabbi of L was Meshulam Zalman Ashkenazi, who was a Mitnaged, but conciliatory towards the Chassidim. On his death the struggle for the rabbinate was renewed - and this time too the mitnagdim carried the day. The new Rabbi chosen was Dov Berish Ashkenazi from the faction "Noda Beshearim" (Known in the Gates). He followed in his predecessor's path, and enjoyed general support in the town. During his term of office the new "Parnas" Bet Midrash was built and its activities recorded. Upon his death in 1852, his son, Yaakov Heshel Ashkenazi was chosen as Rabbi, and he fought a bitter war against Chassidism, which had once more established itself in L with the arrival of R. Yehuda Leib Eiger, who established a new dynasty in the town. This dynasty continued until the Holocaust through his son, R. Avraham Eiger, and his grandsons, R. Azriel Meir and R. Shlomo. The last two perished under the German occupation, the former in 1940 and the latter in 1941.
In 1868 the Chassidic Rabbi, Shneur Zalman Ladjer, author of "Torat Chesed" (Torah of Mercy), was appointed Rabbi of L. In 1892 he left for Palestine, and was replaced by a Mitnaged - R. Hillel Aryeh Lifszyc, who maintained good relations with the Chassidim. Two Chassidic candidates stood against him - R. Gershon Chanoch of Radzin and R. Meir Yechiel of Ostrowiec, but R. Lifszyc won; he was also a man with a more liberal educational background.
During this period several Admorim settled in L, among them R. Tsadok Hacohen, who inherited some of the disciples of R. Yehuda Leib Eiger and was renowned as a great Jewish thinker; R. Moshe Mordechai Twerski, author of "Maamar Mordechai" (The Sayings of Mordechai); and many others.
In 1910 R. Eliahu Klackin (the father of Professor Yaakov Klackin) was chosen Chief Rabbi. He emigrated to Palestine in 1925 and settled in Jerusalem, where he died in 1932. He wrote several important books, that are mentioned in his last work, "Dvarim Achadim" (A Few Things").
In addition to the community institutions, organisations and synagogues mentioned above, note should also be made of other institutions that arose and operated in the 19th century. In 1886 a Jewish hospital was opened, and at the end of the century it contained 90 beds. The community also had an Old Age Home and an Orphanage, and in almost every Jewish quarter there were various charitable organisations alongside the local prayer-houses.
Towards the end of the century the children of the community - 800 boys and 100 girls - were studying in 43 private cheders. The community also ran a Talmud-Torah for children of the poor. During this period non-Jewish private schools also sprang up in the town, first for girls and afterwards for boys. Their language of instruction was Russian or Polish. They also had Jewish pupils. In the 1890s a government school for Jewish children was established in the Catholic Cathedral building, and the language used here was Russian. At one time it had some 300 pupils; but it closed when the Russians withdrew from the town in 1915. In 1897/8 two elementary schools were opened, with Hebrew as the language of instruction, based on the system "Hebrew through Hebrew". One of them, the Pines School, functioned until 1907; the other until the First World War. 1913 also saw the establishment of the "Yavneh" school, which closed in 1927 for lack of money, and after the headmaster, Szmuel Rotensztajn, had emigrated to Palestine.
In addition to the above there were also government elementary schools for Jewish children, beginning with the Austrian occupation and continuing into the 1930s. They were closed on Sabbaths and Holydays, and were therefore known as "Szabsuvka".
An important many-sided cultural institution was "Hazamir" (The Nightingale), established in 1908 with the aim of fostering a love of Hebrew songs and literature among the young. Hazamir also possessed a library of Yiddish, Hebrew, Russian and Polish books - the first modern Jewish public library in L; and it also boasted an orchestra, and a drama group, that gave performances of Jewish classical plays in L and the nearby provincial towns. The rooms of Hazamir were also the setting for lectures on current and Jewish subjects, with some speakers coming from as far away as Warsaw. In 1910 the organisation "Chovevei Sfat Aver" (Lovers of the Hebrew Tongue) was founded; it held Hebrew courses and published Hebrew books. During the First World War, however, this organisation folded, and its assets, including the library, passed into the hands of the "Bund".
At the end of the 19th century the first groups of "Chovevei Zion" (Lovers of Zion) appeared; to begin with their activities were more social than political. At the beginning of the 20th century small groups of the "Bund" appeared on the L scene; in 1903 its first pamphlets were distributed; and in the following year a branch was officially established. In 1905, the year of revolution, the Bund attracted many members, who organized demonstrations and strikes. Activity also increased of the Jewish groups affiliated to the PPS, the Polish Socialist Party. Folders and booklets in Yiddish were distributed in the streets. A small group of Jewish anarchists also came on the scene. The activities of the Bund resulted in some of its members being exiled to Siberia. A Jewish police agent was murdered, and during his funeral disturbances ensued and the Cossacks were called upon to restore order.
The outbreak of the First World War brought with it clear signs of antisemitism by the Russian rulers, who accused the Jews of treachery and aiding the enemy.The Cossacks of the garrison were permitted to plunder Jewish property and to forego payment for goods acquired from Jewish merchants. Some of the town's Jews were exiled to the depths of Russia. Another affliction that affected many Jewish families was forced military conscription. Many young Jews were mobilized, and soon messages began to arrive of their dear ones being killed or wounded in battle. As the front drew near many refugees streamed to L, the economy of the town stagnated, and many Jews were left without work or livelihood.
In July 1915 the Russians withdrew from the town and gave way to the Austrians. Before the retreat the Russian commander summoned the engineer Tomorowicz, one of the leading citizens, and ordered him to organize a militia to guarantee order in the town. Tomorowicz was a known liberal, and he also saw to it that a militia was formed in the Jewish quarter ("the Sixth Commissariat"). The head of this militia was Bernard Glowinski. It numbered some 200 men aged 18 and above, and who had an unblemished record. Its members were armed with revolvers. A tribunal was also set up with judges from among reputable Jewish lawyers in the town. This militia, which consisted only of volunteers, carried out its duties impeccably and to the satisfaction of the citizens. At the end of November 1915, however, it was disbanded and replaced by professionals, but many of the Jewish militiamen joined the new body.
During the Austrian occupation the existing Jewish political organisations increased their activities, and new ones were founded, including youth movements and groups.
The Jewish community and the various organisations of the time opened public soup kitchens that provided warm meals and food to the needy, and especially to children.
With the cessation of hostilities and the proclamation of the independent state of Poland in November 1918, the Jews of L, and Poland in general, hoped for a better future in the new nation. They set about repairing the ravages of war and looked forward to a normalisation of economic, political and social life.
However, as early as 1919 events took place in the town that cast a shadow over the community. On a Thursday, which was the market day, Polish army recruits ran amok, attacked Jewish passers-by, destroyed and even looted Jewish shops. Jewish youngsters organised self-defence, and to begin with even managed to repulse the rioters; but in the afternoon there burst onto the scene in Lubartowska Street gangs of butchers and simple trouble-makers, with butchers' knives, iron bars, and other weapons, and began a pogrom that went on until evening. The police did not interfere. The leaders of the Jewish community met with the Mayor, and only then were the police ordered to stop the pogrom. The results of these disturbances were 60 injured Jews, of whom three died, and widespread damage to shops and houses. Thus was Jewish L included in the many Jewish communities in Poland that were victims of the riots of 1918-19, after Poland's resurgence.
When the disturbances were over the Jews began to reconstruct their businesses and the institutions of the community. Their first concern was the wretched economic situation following the war. In 1921, after two years of rehabilitation, there was still no renewal of activity in 57 of the 1,714 Jewish-owned workshops and businesses, some of which did not employ hired labour. The units that did begin to operate, 1,657 in number, were mostly small businesses that all in all employed 4,871 persons, half of whom (2,366 - or 48.5%) were the owners and their families. The remainder were employees (2,068 Jews and 437 non-Jews). Most of these businesses made clothing (made -to-measure and ready-made), and hats (to order and for shops); there were also furriers, shoemakers, and leather workshops. Together these employed 2,444 persons, of whom 1,031 were workshop owners, 394 were members of their families, and 997 Jews and a mere 22 non-Jews were hired workers. The next largest group were 139 food businesses, employing 804 persons, more than half of them (424 Jews and 129 non-Jews) employees. The remaining plants and workshops were engaged in a variety of manufactures - timber, building, etc., with a small number of participants compared to the leading segments of clothing and food. Most of the workers in production, particularly the hired ones, were men; a few of the hired workers were children (116 Jews and 34 non-Jews). Most businesses were only active part of the year (construction work lay dormant in the winter, and the manufacture of clothing was also seasonal). In the "dead" periods the unemployed and their would-be employers waited for "better times". In the larger plants there were few Jews employed, partly because they were not welcome, and partly because the Jews themselves did not want to work in them (since work continued on the Sabbath and Jewish Holydays). Jews were not employed in government offices or the municipality, in public transport, or even in cleaning the streets.
The most important sphere after crafts and minor industry were the various forms of trading - wholesale, retail, peddling, and market stallholding. Most of those involved were petty traders, who set up market stalls or peddled their wares in the streets of the town or in the surrounding villages. Their volume of goods and of course their turnover were limited. Their customers were the Polish peasants or the labourers, who were economically poorly off and also subject to unemployment. Most of the peasants had smallholdings that supplied nearly all their needs, though there were articles they had to buy outside, such as tobacco, salt, spirits and matches, all of which were state monopolies. Government policy was to squeeze Jewish traders out of this monopoly sector.
One branch was a kind of Jewish monopoly: the leather industry. In the years preceding the Second World War some 95% of production and tanning of hides was in Jewish hands. In 1939 150 workers, 75 of them Jews, were employed in the large tanning plant in the town. The annual turnover of this branch amounted to 20 million zloty (about 4 million dollars). Under Jewish ownership were also a brandy distillery, a brewery, brickworks, flour mills, and a cigarette and tobacco factory, established in 1860 and at its peak employing 400 workers, all of them Jews. After the war, however, this factory was taken over by the State Monopoly and its Jewish workers dismissed.
Branches employing a majority of Jews were the ones to suffer most in times of crisis. Usually it was not only Jewish apprentices and salesmen who became unemployed, but also the workshop owners, most of whom were self-employed. These latter were not entitled to unemployment benefits, nor were they registered in a sick fund, and they and their families were thus unable to receive medical treatment.
In the early years of the 20s the Jews of L, as in other places, set about restoring the mutual aid societies that had existed before the war, and establishing new ones. In the inter-war period there were in L ten provident funds affiliated to the prayer-houses and the political organisations and trade unions. The largest of them was set up immediately after the war by the "Va'ad Hahatsala" (Rescue Committee), and its equity capital amounted in the 30s to 140,000 zloty, and its annual turnover to more than 11 million zloty. This fund gave loans at a low rate of interest, and sometimes at no rate of interest, mainly to workshop owners and petty merchants. These loans were often the only way in which shopkeepers and pedlars could replenish their stocks or replace workshop equipment.
In this period too ten Jewish banks were established, the first of them, "Kupat Halvaa Uchisachon" (Loan and Savings Bank) at the end of the 19th or the beginning of the 20th century. Later there appeared "Chevra Le'ezra Hadadit Lefkidei Mischar Usocharim Ze'irim Bnei Dat Moshe" (Mutual Aid Society of Trading Clerks and Small Traders of the Mosaic Faith); "Habank Harishon Levaalei Mlacha" (First Bank for Artisans); "Bank Hasocharim" (Merchants' Bank); "Bank Baalei Batim" (Householders' Bank); and "Bank Shitufi" (Cooperative Bank). The employees too established their own mutual aid societies in conjunction with the trade unions. Other bodies were set up at the beginning of the 20th century within the framework of existing organisations. In the 20s there was widespread organized professional activity, embracing most of the Jewish workers. One of the most prominent organisations, in terms of number of members and level of activity was the Union, or Guild, of Needleworkers, founded as early as 1915. Other unions of considerable size followed, such as those of the cobblers, printing workers, tanners, carpenters, butchers, porters, and bakery workers. The largest of them, "Haigud Hameuchad shel Poalei O'r" (The United Guild of Leather Workers) counted some 500 members; while the Textile Workers' Guild numbered some 200 members.
The community of L was blessed with many institutions of welfare, charity and aid. Especially renowned was the society "Achiezer" (Brotherly Aid), which on the threshold of Sabbaths and Festivals distributed food to hundreds of the needy. Money for its activities came from contributions from members of the community (at Purim, on the Eve of Yom Kippur and at weddings). The institution also ran a kitchen for the poor. The organisation "Linat Tsedek" (Hospice for the Poor) provided the sick and destitute with medical care and drugs free of charge, and its members helped their families stand watch by the bedside of the ill. In 1919 the organisation "Zichron Nes" ( Token of Miracle) was inaugurated with the aim of helping war invalids, bereaved parents, and orphans. It was financed from the sale of the privilege of being called to the Torah ("Kearot Nedava" - Collections of Charity) on the occasion of weddings etc. Members of the "Shabta Tava" (Good Sabbath) society were wont to walk through the streets on a Friday to collect chalot (Sabbath bread) for the needy. The society "Hachnasat Orchim" (Hospice) opened a lodging for indigent travellers: each night tens of them were accommodated. In addition to all these bodies other groups of community volunteers were active, for example, in visiting the vick, supporting the poor, collecting a dowry for penniless brides, and other causes.
In 1923 a branch of "TOZ" (Polish initials for the Jewish Health Organisation of Poland) was established in L. Its activities were many and varied: in its clinics there were 22 doctors and X-ray and radiography facilities. TOZ ran an advice service for women and a child care centre, and even supported Jewish sports organisations and provided them with training areas. Each summer TOZ organised camps for school pupils. In 1924, the first year these took place, 300 children participated; in 1939 this number reached 1,200.
Another institution the community boasted was the Jewish Hospital, opened, as stated, in 1886, but not allowed to function during the Great War. Afterwards, however, it quickly resumed and expanded its activity, with the support of former residents of L in the USA. It contained 30 beds for men and 26 for women, and had an X-ray department, a bacteriological laboratory, and a clinic for quartz radiation. In the late 30s a children's department was opened, and became known for its high professional standard; and an Institute for Hydrotherapy was also inaugurated. Every week the hospital held lectures and symposia, with the participation of doctors from L and the vicinity.
The inter-war years saw increased activity by the Jewish political parties and youth movements. Prominent among the parties were the various Zionist factions. Zionist activities in L began, as stated, with the groups of "Chovevei Zion" that sprang up at the end of the preceding century. The Zionist Federation crystallised in L during the Austrian occupation (1915-18), and its activity increased following the Balfour Declaration in 1917. When the Mandate for Palestine was decided upon there was much celebration in the community, and all the Jewish shops were closed. At the same time the young people set up "Tseirei Zion" (Youth of Zion), which existed until it split into "Socialist Zionist Youth" and "Poalei Zion Smol" (Leftist Workers of Zion). In 1920 a branch of "Hechalutz" (The Pioneer) was set up. In the Zionist Federation building in L lectures and evening classes were held. The Zionists in L were active in all speheres of the political and social life, took part in elections to the Polish parliament, were represented on the Town Council, and of course in the Community Council.
1924 saw the establishment of "Haliga Lema'an Eretz-Izrael Haovedet" (The League for Labour Israel), with the participation of Socialist Poalei Zion, Hashomer Hatsair, Hechalutz Hatsair, Dror (Freedom), and Hapoel. In 1929 the organisation "Haoved" (The Worker) was launched in L, consisting of the workshop owners, and particularly those who intended to emigrate to Palestine. After the schism in "Poalei Zion" in 1921 the majority of the members of the L branch remained loyal to the left wing of the party. This party exercised much influence among students and pupils; its youth movement, "Jugent", had some 140 members in L. In the late 30s, when the atmosphere in the whole of Poland was permeated with national-fascist ideas, the party was persecuted by the authorities, who regarded members of Left Poalei Zion as clandestine communits, and in 1936-37 some of them were even arrested. Prominent among the Zionist youth movements was "Hashomer Hatsair" (The Young Watchman). The L cell was set up in 1916, and between the two world wars its membership and activity increased, and a kibbutz training farm was even established. There was also in L a large branch of "Hanoar Hazioni", the youth movement of the General Zionists. In 1927 was founded the children's organisation in the name of Borochov, called "Jungbor", which was active in the fields of education, information and scouting.
The L branch of "Hamizrachi" was founded in 1903, and was reorganized at the end of the Great War. Its activities were centred on the Mizrachi Synagogue in Novorowna Street, while for its lectures and discussions it rented the "Hazamir" hall once a week. Its members opened the "Yavneh School" in L, which turned out several graduation classes. Hamizrachi's youth movement (Tseirei Mizrachi) established several training centres outside the town. In the early 20s branches of "Hapoel Hamizrachi" and "Hechalutz Hamizrachi" were also set up.
The Revisionist Movement was active in L from 1925 on. The first branch to be set up was "Hashachar" (The Dawn), and a year later saw the arrival of the youth movement "Beitar" (Covenant of Josef Trumpeldor). The Revisionist movement in L and district expanded, and in 1934 had some 50 branches. In the same period one of its important fractions, called "Brit Hachayil" (Covenant of Strength), also developed. Beitar ran several amateur groups - drama, literature, etc., as well as a brass band. Each year it organised summer camps, where the young boys and girls received instruction with regard to emigrating to Palestine. To the congresses of the L branch of the Revisionists there came speakers from the national Polish and the world organisation (including Ze'ev Jabotinsky and Menachem Begin) and their meetings attracted a large audience.
Nevertheless, not all the Jews of L belonged to the Zionist camp. Up to 1914 the "Bund" operated clandestinely, but renewed its public activities during the Austrian occupation of 1915-18. During the war there was a shortage of food bordering on famine, and the members of the Bund therefore opened a kitchen and set up cooperatives of consumers and a bakery, which employed ten workers. Bund members were also active in the field of culture, and among other things established a library (in the name of Groser), a choir, and a drama group. In 1917 was held the L Congress, where it was decided to sever all connection with the Bund of the Russian Empire and to maintain an independent Polish section. In November 1917, after the October Revolution in Russia, a Workers' Council was set up in L, also with the participation of the Bund. At this time there were attacks on the Jews in L, and to counter them a militia of Polish and Jewish members of the Polish Socialist Party was formed. In 1920 during the Polish-Soviet War the Polish authorities arrested some of the active leaders of the Bund in L. After the war the Bund played an influential role in the Central Committee of the Jewish trade unions in L, with five representatives; and in most of the individual unions it formed a majority, while Poalei Zion and the communists usually had only one representative. The youth and children's movements of the Bund in L - "Zukunft" (The Future) and "Sakif" ( Threshold ?? ) - had large and active branches.
Some L Jews were members of the Communist Party, which operated several illegal cells in the town. These members often infiltrated the other Jewish workers' parties (Poalei Zion and Bund). They were mainly active in organising strikes and demonstration and in spreading propaganda. During the 30s some active Jewish communists were arrested and imprisoned; and the threat of imprisonment impelled others to flee abroad, particularly to France, where they continued to maintain contact and help one another.
Also present in L were "Agudat Israel", its labour faction "Poalei Agudat Israel", and its youth movement "Tseirei Agudat Israel". Most members of Aguda were to be found among the various Chassidic groups; and their field of action was within the Community Council and in religious education.
Up to 1932 the majority in the community council were religious representatives, and only a few came from among the "Neurim" (The Enlightened or Illuminati), as secular Jews were then called. The council consisted of the general body, or council, and the executive committee. In the inter-war period elections to the council were only partly democratic; women were not eligible to vote, and there were also restrictions with regard to age and economic status. From 1932 on the elections were held on the basis of party lists, and in that year all the parties that had branches in L participated, with the exception of the Bund, which boycotted the elections. Five representatives of Agudat Israel and the Chassidim were elected, together with two from the Artisan List, and one each from the General Zionists, Socialist Poalei Zion, Left Poalei Zion, Mizrachi, the "People's Party", and the "Assimilationists". At first the Chairman chosen was the lawyer M. Alten, but the Polish authorities withheld their approval, and his place was taken by Y. Zilber. As Chairman of the executive committee, the "assimilationist", the lawyer A. Levinsohn, was chosen. He held this post for only a year, when he resigned, and his place was taken by S, Halberstadt from Agudat Israel.
In 1930 R. Meir Shapira (1886-1933), one of the leaders of Agudat Israel, was elected Chief Rabbi of L. He established the largest Yeshiva in Poland, "Yeshivat Chachmei Lublin" (The Yeshiva of the Wise Men of L). R. Shapira was a brilliant speaker, and was a member of the Polish parliament (Sejm). He initiated learning of the ”Daf Hayomi" (The Daily Page) throughout the Diaspora. On his death there were no more chief rabbis in L, and the community contented itself with dayanim (judges).
The closing years of the 20s and the beginning of the 30s found the community in grave circumstances. It was weighed down with heavy debts, and its income, mainly accruing from schechita (ritual slaughtering), payment for burial plots, and direct taxes, did not cover expenditure, viz., wages in the religious sector (in the Rabbinate there were a rabbi, five dayanim, and a number of "minor" rabbis in the suburbs); allocations to various religious institutions and a Talmud-Torah, to the Yeshivat Chochmei Lublin, to the orphanage, to schools where the language was Yiddish; and so forth. Due to the economic problems these allocations were not made on time. and the institutions of the community therefore operated at a low level. The orphanage, for example, was almost without supervision, the children were neglected, and some of them became subject to bad influences, The cemetery too was neglected, and sometimes waterlogged, and there was no vacant ground for new burials.
In 1932, following election to the new council, its members set about a thorough reorganisation. Hours of work and reception times for the public were scheduled; the collection of taxes came under strict control; renovations were carried out on the synagogues, especially that of the Maharashal; and the amount and times of payment for allocations to the various institutions were laid down.
In the elections of 1939 members of the Bund were also represented. The community faced a crisis, when the authorities insisted on the shochtim (ritual slaughterers) being given functionary status. These schochtim exploited the situation by refusing to pay slaughtering fees to the community, and this seriously undermined its finances. One of the last decisions of the community before the Holocaust was that of August 7th, 1939 - to grant the Jewish Hospital an allocation of 6,000 zloty, and to issue bonds to the samount of 30,000 zloty to ensure normal operation of the council and its subsidised institutions.
The Jews of L were represented in the town council via the party lists. In the elections of 1927 16 Jews (most of them from the Bund) were elected to the council, i.e. a third of all its members. At its meetings they put forward the claims of the Jewish institutions and of the Jews of L in general, and protested against all forms of discrimination.
From earliest times there were in L cheders, where Chumesh, Mishna and Talmud were studied. Poor children went to the Talmud Torah, subsidized by the community. Older boys studied at Batei Midrash; and in the evenings, by candlelight, lessons were held in Gemara, Mishna, or the weekly passage of the Torah for adults and married men - and there were among these some veritable scholars.
In the inter-war period L, as in other parts of Poland, possessed some state elementary schools for Jewish children - the Szabasuvka. After Polish independence three additional schools of this type were opened.
* For explanations of these and other organisations and movements mentioned in this survey readers are kindly referred to more detailed reference works. Back
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CONTRIBUTION OF MUSLIMS SCIENTIST IN THE FIELD OF MEDICINE AND PHARMACY
In the Modern Muslim world is seen as many things , but rarely is it viewed as a source of inspiration and enlightenment . Though it is a force of enlightenment and it is not only verses of the Quran that testy to that fact, but also the great body scholarship produced during the Middle Ages. While Europe was in the midst of darkness, it was the Muslims spurred on the light of their news who picked up the torch of scholarship and science .It was the Muslims who preserved the knowledge of old days , elaborate upon it and passed it to Europe .
Islamic scientist in Medicine and Pharmacy was not known for many age due to many factors including a) Colonialism and lack of knowledge among Muslims generation in this modern world, other factors is b) lack of institution based facts how and when Islamic Scientists what contribution was . Although every people earn what they do pass on it generation after generation. So we the Islam of today learners should make important to look back over and make research topic selected on this issue and gain knowledge of, respected and Prized the contribution Islamic Civilization by early Muslim scholars
The first Muslim physician is believed to have been Muhammad himself, as a significant number of hadiths concerning medicine are attributed to him. Several Sahaba are said to have been successfully treated of certain diseases by following the medical advice of Muhammad. The three methods of healing known to have been mentioned by him were Honey, Cupping, and Cauterization, though he was generally opposed to the use of cauterization unless it “suits the ailment.” According to Ibn Hajar al-Asqalani, Muhammad disliked this method due to it causing “pain and menace to a patient” since there was no anesthesia in his time. Muhammad also appears to have been the first to suggest the contagious nature of leprosy, mange and sexually transmitted disease; and that there is always a cause and a cure for every disease, according to several hadiths in the Sahih al-Bukhari, Sunan Abi Dawood and Al-Muwatta attributed to Muhammad, such as:
“There is no disease that Allah has created, except that He also has created its treatment.
“Make use of medical treatment, for Allah has not made a disease without appointing a remedy for it, with the exception of one disease, namely old age.
“Allah has sent down both the disease and the cure, and He has appointed a cure for every disease, so treat you medically.
“The one who sent down the disease sent down the remedy.
The belief that there is a cure for every disease encouraged early Muslims to engage in biomedical research and seek out a cure for every disease known to them. Many early authors of Islamic medicine, however, were usually clerics rather than physicians, and were known to have advocated the traditional medical practices of prophet Muhammad’s time, such as those mentioned in the Qur’a n and Hadith. For instance, therapy did not require a patient to undergo any surgical procedures at the time.
From the 9th century, Hunayn ibn Ishaq translated a number of Galen‘s works into the Arabic language, followed by translations of the Sushruta Samhita, Charaka Samhita and Middle Persian works from Gundishapur. Muslim physicians soon began making many of their own significant advances and contributions to medicine, including the fields of
ALLERGOLOGY, ANATOMY, BACTERIOLOGY, BOTANY, DENTISTRY, EMBRYOLOGY, ENVIRONMENTALISM, ETIOLOGY, IMMUNOLOGY, MICROBIOLOGY, OBSTETRICS, OPHTHALMOLOGY, PATHOLOGY, PEDIATRICS, PERINATOLOGY, PHYSIOLOGY, PSYCHIATRY, PSYCHOLOGY, PULSOLOGY AND SPHYGMOLOGY, SURGERY, THERAPY, UROLOGY, ZOOLOGY, AND THE PHARMACEUTICAL SCIENCES SUCH AS PHARMACY AND PHARMACOLOGY, AMONG OTHERS.
Medicine was a central part of medieval Islamic culture. Responding to circumstances of time and place, Islamic physicians and scholars developed a large and complex medical literature exploring and synthesizing the theory and practice of medicine Islamic medicine was initially built on tradition, chiefly the theoretical and practical knowledge developed in Arabia, Persia, Greece, Rome, and India. Galen and Hippocrates were pre-eminent authorities, as well as the Indian physicians Sushruta and Charaka, and the Hellenistic scholars in Alexandria. Islamic scholars translated their voluminous writings from Greek and Sanskrit into Arabic and then produced new medical knowledge based on those texts. In order to make the Greek and Indian traditions more accessible, understandable, and teachable, Islamic scholars ordered and made more systematic the vast and sometimes inconsistent Greco-Roman and Indian medical knowledge by writing encyclopedias and summaries. It was through Arabic translations that the West learned of Hellenic medicine, including the works of Galen and Hippocrates. Of equal if not of greater influence in Western Europe were systematic and comprehensive works such as Avicenna‘s The Canon of Medicine, which were translated into Latin and then disseminated in manuscript and printed form throughout Europe. During the fifteenth and sixteenth centuries alone, The Canon of Medicine was published more than thirty-five times.
Muslim physicians set up the earliest dedicated hospitals in the modern sense, known as Bimaristans, which were establishments where the ill were welcomed and cared for by qualified staff, and which were clearly distinguished from the ancient healing temples, sleep temples, hospices, assylums, lazarets and leper-houses which were more concerned with isolating the sick and the mad from society “rather than to offer them any way to a true cure. The Bimaristan hospitals later functioned as the first public hospitals, psychiatric hospitals and diploma-granting medical universities.
In the medieval Islamic world, hospitals were built in all major cities; in Cairo for example, the Qalawun Hospital could care for 8,000 patients, and a staff that included physicians, pharmacists, and nurses. One could also access a dispensary, and research facility that led to advances, which included the discovery of the contagious nature of diseases, and research into optics and the mechanisms of the eye. Muslim doctors were removing cataracts with hollow needles over 1000 years before Western physicians dared attempt such a task. Hospitals were built not only for the physically sick, but for the mentally sick also. One of the first ever psychiatric hospitals that cared for the mentally ill was built in Cairo. Hospitals later spread to Europe during the Crusades, inspired by the hospitals in the Middle East. The first hospital in Paris, Les Quinze-vingts, was founded by Louis IX after his return from the Crusade between 1254-1260.
Hospitals in the Islamic world featured competency tests for doctors, drug purity regulations, nurses and interns, and advanced surgical procedures. Hospitals were also created with separate wards for specific illnesses, so that people with contagious diseases could be kept away from other patients.
One of the features in medieval Muslim hospitals that distinguished them from their contemporaries and predecessors was their significantly higher standards of medical ethics. Hospitals in the Islamic world treated patients of all religions, ethnicities, and backgrounds, while the hospitals themselves often employed staff from Christian, Jewish and other minority backgrounds. Muslim doctors and physicians were expected to have obligations towards their patients, regardless of their wealth or backgrounds. The ethical standards of Muslim physicians was first laid down in the 9th century by Ishaq bin Ali Rahawi, who wrote the Adab al-Tabib (Conduct of a Physician), the first treatise dedicated to medical ethics. He regarded physicians as “guardians of souls and bodies”, and wrote twenty chapters on various topics related to medical ethics.
Another unique feature of medieval Muslim hospitals was the role of female staff, who were rarely employed in ancient and medieval healing temples elsewhere in the world. Medieval Muslim hospitals commonly employed female nurses, including nurses from as far as Sudan, a sign of great breakthrough. Muslim hospitals were also the first to employ female physicians, the most famous being two female physicians from the Banu Zuhr family who served the Almohad ruler Abu Yusuf Ya’qub al-Mansur in the 12th century. Later in the 15th century, female surgeons were illustrated for the first time in Şerafeddin Sabuncuoğlu‘s Cerrahiyyetu’l-Haniyye (Imperial Surgery)
The first encyclopedia of medicine in Arabic was Ali ibn Sahl Rabban al-Tabari‘s Firdous al-Hikmah (“Paradise of Wisdom”), written in seven parts, c. 860. It was the first to deal with pediatrics and child development, as well as psychology and psychotherapy. In the fields of medicine and psychotherapy, the work was primarily influenced by Islamic thought and ancient Indian physicians such as Sushruta and Charaka. Unlike earlier physicians, however, al-Tabari emphasized strong ties between psychology and medicine, and the need of psychotherapy and counseling in the therapeutic treatment of patients
Muhammad ibn Zakarīya Rāzi (Rhazes) wrote the Comprehensive Book of Medicine in the 9th century. The Large Comprehensive was the most sought after of all his compositions, in which Rhazes recorded clinical cases of his own experience and provided very useful recordings of various diseases. The Comprehensive Book of Medicine, with its introduction of measles and smallpox, was very influential in Europe.
(HALY ABBAS)’S KITAB KAMIL AS-SINA’A AT-TIBBIYYA (“Complete Book of the Medical Art“), c. 980, became better known as the Kitab al-Maliki (“Royal Book“, Latin: Liber regalis) in honour of its royal patron ‘Adud al-Dawla. In twenty sections, ten of theory and ten of practice, it was more systematic and concise than Razi’s Hawi, but more practical than Avicenna’s Canon, by which it was superseded. With many interpolations and substitutions, it served as the basis for the Pantegni (c. 1087) of Constantinus Africanus, the founding text of the Schola Medica Salernitana in Salerno
(Abulcasis), regarded as the father of modern surgery contributed greatly to the discipline of medical surgery with his Kitab al-Tasrif (“Book of Concessions“), a 30-volume medical encyclopedia published in 1000, which was later translated to Latin and used in European medical schools for centuries. He invented numerous surgical instruments and described them in his al-Tasrif.
Avicenna (Ibn Sina), a Hanbali and Mu’tazili philosopher and doctor in the early 11th century, was another influential figure. He is regarded as the father of modern medicine, and one of the greatest thinkers and medical scholars in history. His medical encyclopedia, The Canon of Medicine (c. 1020), remained a standard textbook in Europe for centuries, up until the renewal of the Muslim tradition of scientific medicine. He also wrote The Book of Healing (actually a more general encyclopedia of science and philosophy), which became another popular textbook in Europe. Among other things, Avicenna’s contributions to medicine include the introduction of systematic experimentation and quantification into the study of physiology,the discovery of the contagious nature of infectious diseases, the introduction of quarantine to limit the spread of contagious diseases, the introduction of experimental medicine, evidence-based medicine, clinical trials,randomized controlled trialsefficacy tests, clinical pharmacology, risk factor analysis, and the idea of a syndrome in the diagnosis of specific diseases,the first descriptions on bacterial and viral organisms, the distinction of mediastinitis from pleurisy, the contagious nature of phthisis and tuberculosis, the distribution of diseases by water and soil, and the first careful descriptions of skin troubles, sexually transmitted diseases, perversions, and nervous ailments, as well the use of ice to treat fevers, and the separation of medicine from pharmacology, which was important to the development of the pharmaceutical sciences.
Ibn al-Nafis (1213-1288) wrote Al-Shamil fi al-Tibb (The Comprehensive Book on Medicine), a voluminous medical encyclopedia that was originally planned to comprise 300 volumes, but he was only able to complete 80 volumes as a result of his death in 1288. However, even in its incomplete state, the book is one of the largest known medical encyclopedias in history, though only a small portion of The Comprehensive Book on Medicine has survived. During his lifetime, The Comprehensive Book on Medicine had eventually replaced Ibn Sina’s The Canon of Medicine as a medical authority in the medieval Islamic world. Arabic biographers from the 13th onwards considered Ibn al-Nafis the greatest physician in history, some referring to him as “the second Ibn Sina”, and others considering him even greater than Ibn Sina
In the 9th century
AL-KINDI (ALKINDUS), IN DE GRADIBUS Demonstrated the application of mathematics and quantification to medicine, particularly in the field of pharmacology. This includes the development of a mathematical scale to quantify the strength of drugs, and a system that would allow a doctor to determine in advance the most critical days of a patient’s illness, based on the phases of the Moon.
In the 10th century, RAZI (RHAZES) Introduced controlled experiment and clinical observation into the field of medicine, and rejected several Galenic medical theories unverified by experimentation.The earliest known medical experiment was carried out by Razi in order to find the most hygienic place to build a hospital. He hung pieces of meat in places throughout 10th century Baghdad and observed where the meat decomposed least quickly, and that was where he built the hospital. In his Comprehensive Book of Medicine, Razi recorded clinical cases of his own experience and provided very useful recordings of various diseases. In his Doubts about Galen, Razi was also the first to prove both Galen‘s theory of humorism and Aristotle‘s theory of classical elements false using experimentation He also introduced urinalysis and stool tests
AVICENNA (IBN SINA)
Is considered the father of modern medicine,for his introduction of systematic experimentation and quantification into the study of physiology, the introduction of experimental medicine, clinical trials, risk factor analysis, and the idea of a syndrome in the diagnosis of specific diseases, in his medical encyclopedia, The Canon of Medicine (c. 1020), which was also the first book dealing with evidence-based medicine, randomized controlled trials,and efficacy tests.
According to Toby Huff and A. C. Crombie, the Canon contained “a set of rules that laid down the conditions for the experimental use and testing of drugs” which were “a precise guide for practical experimentation” in the process of “discovering and proving the effectiveness of medical substances. The Canon laid out the following rules and principles for testing the effectiveness of new drugs and medications, which still form the basis of clinical pharmacologyand modern clinical trials:
- “The drug must be free from any extraneous accidental quality.”
- “It must be used on a simple, not a composite, disease.”
- “The drug must be tested with two contrary types of diseases, because sometimes a drug cures one disease by Its essential qualities and another by its accidental ones.”
- “The quality of the drug must correspond to the strength of the disease. For example, there are some drugs whose heat is less than the coldness of certain diseases, so that they would have no effect on them.”
- “The time of action must be observed, so that essence and accident are not confused.”
- “The effect of the drug must be seen to occur constantly or in many cases, for if this did not happen, it was an accidental effect.”
- “The experimentation must be done with the human body, for testing a drug on a lion or a horse might not prove anything about its effect on man.”
IBN ZUHR (AVENZOAR)
Who introduced the experimental method into surgery, for which he is considered the father of experimental surgery.Other early supporters of human dissection and autopsy include Ibn Tufail,Saladin‘s physician Ibn Jumay, Abd-el-latif,and Ibn al-Nafis
The experimental method was introduced into botany, materia medica and the agricultural sciences in the 13th century by the Andalusian-Arab botanist Abu al-Abbas al-Nabati, the teacher of Ibn al-Baitar. Al-Nabati introduced empirical techniques in the testing, description and identification of numerous material medica, and he separated unverified reports from those supported by actual tests and observations.
ANATOMY AND PHYSIOLOGY: In anatomy and physiology, the first physician to refute Galen‘s theory of humorism was
MUHAMMAD IBN ZAKARĪYA RĀZI (RHAZES) In his Doubts about Galen in the 10th century. He criticized Galen’s theory that the body possessed four separate “humors” (liquid substances), whose balance are the key to health and a natural body-temperature. Razi was the first to prove this theory wrong using an experiment. He carried out an experiment which would upset this system by inserting a liquid with a different temperature into the body resulting in an increase or decrease of bodily heat, which resembled the temperature of that particular fluid. Razi noted particularly that a warm drink would heat up the body to a degree much higher than its own natural temperature, thus the drink would trigger a response from the body, rather than transferring only its own warmth or coldness to it. This line of criticism was the first comprehensive experimental refutation of Galen’s theory of humours and Aristotle‘s theory of the four classical elements on which it was grounded. Razi’s own chemical experiments suggested other qualities of matter, such as “oiliness” and “sulfurousness“, or inflammability and salinity, which were not readily explained by the traditional fire, water, earth and air division of elements.
EXPERIMENTAL ANATOMY AND PHYSIOLOGY
The contributions of Avicenna to physiology include the introduction of systematic experimentation and quantification into the study of physiology in The Canon of Medicine (c. 1020). The contributions of Ibn al-Haytham (Alhacen) to anatomy and physiology include his correct explanation of the process of sight and visual perception for the first time in his Book of Optics, Other innovations introduced by Muslim physicians to the field of physiology by this time include the use of animal testing.
Ibn Zuhr (Avenzoar) (1091-1161) was the first physician known to have carried out human dissections and postmortem autopsy. He proved that the skin disease scabies was caused by a parasite, a discovery which upset the theory of humorism supported by Hippocrates and Galen. The removal of the parasite from the patient’s body did not involve purging, bleeding, or any other traditional treatments associated with the four humours.
In the 12th century, Saladin‘s physician Ibn Jumay was also one the first to undertake human dissections, and he made an explicit appeal for other physicians to do so as well. During a famine in Egypt in 1200, Abd-el-latif observed and examined a large number of skeletons, and he discovered that Galen was incorrect regarding the formation of the bones of the lower jaw and sacrum.
Circulatory anatomy and physiology:
IBN AL-NAFIS The father of circulatory physiology was another early proponent of human dissection In 1242, he was the first to describe the pulmonary circulation, coronary circulation and capillary circulation, which form the basis of the circulatory system, for which he is considered the one of the greatest physiologists in history The first European descriptions of the pulmonary circulation came several centuries later, by Michael Servetus in 1553 and William Harvey in 1628. Ibn al-Nafis also described the earliest concept of metabolism, and developed new Nafisian systems of anatomy, physiology and psychology to replace the Avicennian and Galenic doctrines, while discrediting many of their erroneous theories on the four humours, pulsation, bones, muscles, intestines, sensory organs, bilious canals, esophagus, stomach, and the anatomy of almost every other part of the human body.
The Arab physician Ibn al-Lubudi (1210-1267), also from Damascus, wrote the Collection of discussions relative to fifty psychological and medical questions, in which he rejects the theory of four humours supported by Galen and Hippocrates, discovers that the body and its preservation depend exclusively upon blood, rejects Galen’s idea that women can produce sperm, and discovers that the movement of arteries are not dependant upon the movement of the heart, that the heart is the first organ to form in a fetus‘ body (rather than the brain as claimed by Hippocrates), and that the bones forming the skull can grow into tumors. He also advises that in cases of extreme fever, a patient should not be released from hospital.
PULSOLOGY AND SPHYGMOLOGY: Muslim physicians were pioneers in pulsology and sphygmology. In ancient times, Galen as well as Chinese physicians erroneously believed that there was a unique type of pulse for every organ of the body and for every disease. Galen also erroneously believed that “every part of an artery pulsates simultaneously” and that the motion of the pulse was due to natural motions (the arteries expanding and contracting naturally) as opposed to foced motions (the heart causing the arteries to either expand or contract). The first correct explanations of pulsation were given by Muslim physicians.
“Every beat of the pulse comprises two movements and two pauses. Thus, expansion: pause: contraction: pause. The pulse is a movement in the heart and arteries … which takes the form of alternate expansion and contraction.”
Avicenna also pioneered the modern approach of examining the pulse through the examination of the wrist, which is still practiced in modern times. His reasons for choosing the wrist as the ideal location is due to it being easily available and the patient not needing to be distressed at the exposure of his/her body. The Latin translation of his Canon also laid the foundations for the later invention of the sphygmograph.
Ibn al-Nafis, in his Commentary on Anatomy in Avicenna’s Canon, completely rejected the Galenic theory of pulsation after his discovery of the pulmonary circulation. He developed his own Nafisian theory of pulsation after discovering that pulsation is a result of both natural and forced motions, and that the “forced motion must be the contraction of the arteries caused by the expansion of the heart, and the natural motion must be the expansion of the arteries.” He notes that the “arteries and the heart do not expand and contract at the same time, but rather the one contracts while the other expands” and vice versa. He also recognized that the purpose of the pulse is to help disperse the blood from the heart to the rest of the body. Ibn al-Nafis briefly summarizes his new theory of pulsation:
“The primary purpose of the expansion and contraction of the heart is to absorb the cool air and expel the wastes of the spirit and the warm air; however, the ventricle of the heart is wide. Moreover, when it expands it is not possible for it to absorb air until it is full, for that would then ruin the temperament of the spirit, its substance and texture, as well as the temperament of the heart. Thus, the heart is necessarily forced to complete its fill by absorbing the spirit.”
EPIDEMIOLOGY, ETIOLOGY, PATHOLOGY: In etiology and epidemiology, Muslim physicians were responsible for the discovery of infectious disease and the immune system, advances in pathology, and early hypotheses related to bacteriology and microbiology. Their discovery of contagious disease in particular is considered revolutionary and is one of the most important discoveries in medicine. The earliest ideas on contagion can be traced back to several hadiths attributed to Muhammad in the 7th century, who is said to have understood the contagious nature of leprosy, mange, and sexually transmitted disease. These early ideas on contagion arose from the generally sympathetic attitude of Muslim physicians towards lepers (who were often seen in a negative light in other ancient and medieval societies) which can be traced back through hadiths attributed to Muhammad and to the following advice given in the Qur’an
“There is no fault in the blind, and there is no fault in the lame, and there is no fault in the sick.”
This eventually led to the theory of contagious disease, which was fully understood by Avicenna in the 11th century. By then, the pathology of contagion had been fully understood, and as a result, hospitals were created with separate wards for specific illnesses, so that people with contagious diseases could be kept away from other patients who do not have any contagious diseases. In The Canon of Medicine (1020), Avicenna discovered the contagious nature of infectious diseases such as phthisis and tuberculosis, the distribution of diseases by water and soil, and fully understood the contagious nature of sexually transmitted diseases. In epidemiology, he introduced the method of quarantine as a means of limiting the spread of contagious diseases, and introduced the method of risk factor analysis and the idea of a syndrome in the diagnosis of specific diseases.
In order to find the most hygienic place to build a hospital, Muhammad ibn Zakarīya Rāzi (Rhazes) carried out an experiment where he hung pieces of meat in places throughout 10th century Baghdad and observed where the meat decomposed least quickly. Razi also wrote the Comprehensive Book of Medicine in the 9th century. The Large Comprehensive was the most sought after of all his compositions, in which Razi recorded clinical cases of his own experience and provided very useful recordings of various diseases, as well as the discovery of measles and smallpox. The Large Comprehensive also criticized the views of Galen, after Razi had observed many clinical cases which did not follow Galen’s descriptions of fevers. For example, he stated that Galen’s descriptions of urinary ailments were inaccurate as he had only seen three cases, while Razi had studied hundreds of such cases in hospitals of Baghdad and Rayy.The Comprehensive Book of Medicine, especially with its introduction of measles and smallpox, was very influential in Europe.
IBN ZUHR (AVENZOAR Was the first physician to provide a real scientific etiology for the inflammatory diseases of the ear, and the first to clearly discuss the causes of stridor. Through his dissections, he was also able to prove that the skin disease scabies was caused by a parasite, a discovery which upset the theory of humorism supported by Hippocrates, Galen and Avicenna. He also gave the first accurate descriptions on neurological diseases, including meningitis, intracranial thrombophlebitis, and mediastinal germ cell tumors. Averroes suggested the existence of Parkinson’s disease and attributed photoreceptor properties to the retina. Maimonides wrote about neuropsychiatric disorders and described rabies and belladonna intoxication.
ALLERGOLOGY AND IMMUNOLOGY: The study of allergology and immunology originate from the Islamic world.
MUHAMMAD IBN ZAKARĪYA RĀZI (RHAZES) Was responsible for discovering “allergic asthma“, and was the first physician known to have written articles on allergy and the immune system. In the Sense of Smelling, he explains the occurrence of rhinitis after smelling a rose during the Spring. In the Article on the Reason Why Abou Zayd Balkhi Suffers from Rhinitis When Smelling Roses in Spring, he dicusses seasonal rhinitis, which is the same as allergic asthma or hay fever. Al-Razi was the first to realize that fever is a natural defense mechanism, the body’s way of fighting disease.
The distinction between smallpox and measles also dates back to al-Razi. The medical procedure of inoculation was practiced in the medieval Islamic world in order to treat smallpox. This was later followed by the first smallpox vaccine in the form of cowpox, invented in Turkey in the early 18th century.
DR. ABDIRIZAK HAJI MOHAMED
Pharm-D University of Karachi-Pakistan
Contribution of Muslim scientist is far-reaching topic
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Desalination: A National Perspective 8 A Strategic Research Agenda for Desalination As noted in Chapter 3, desalination is likely to have a niche in the water management portfolio of the future, although the significance of this niche cannot be definitively determined at this time. The potential for desalination to meet anticipated water demands in the United States is not constrained by the source water resources or the capabilities of current technology, but instead it is constrained by financial, social, and environmental factors. Over the past 50 years the state of desalination technology has advanced substantially, and improvements in energy recovery and declining membrane material costs have made brackish water and seawater desalination a more reasonable option for some communities. However, desalination remains a higher-cost alternative for water supply in many communities, and concerns about potential environmental impacts continue to limit the application of desalination technology in the United States. For inland desalination facilities, there are few, if any, cost-effective environmentally sustainable concentrate management technologies. Meanwhile, as noted in Chapter 2, there is no integrated and strategic direction to current federal desalination research and development efforts to help address these concerns. In this chapter, long-term research goals are outlined for advancing desalination technology and improving the ability of desalination to address U.S. water supply needs. A strategic national research agenda is then presented to address these goals. This research agenda is broadly conceived and includes research that could be appropriately funded and conducted in either the public or private sectors. The committee recognizes that research cannot address all barriers to increased application of desalination technology in regions facing water scarcity concerns; therefore, practical implementation issues are discussed separately in Chapter 7. Recommendations related to implementing the proposed research agenda are also provided in this chapter.
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Desalination: A National Perspective LONG-TERM RESEARCH GOALS Based on the committee’s analyses of the state of desalination technology, potential environmental impacts, desalination costs, and implementation issues in the United States (see Chapters 4-7), the committee developed two overarching long-term goals for further research in desalination: Understand the environmental impacts of desalination and develop approaches to minimize these impacts relative to other water supply alternatives, and Develop approaches to lower the financial costs of desalination so that it is an attractive option relative to other alternatives in locations where traditional sources of water are inadequate. Understanding the potential environmental impacts of desalination in both inland and coastal communities and developing approaches to mitigate these impacts relative to other alternatives are essential to the future of desalination in the United States. The environmental impacts of both source water intakes and concentrate discharge remain poorly understood. Although the impacts of coastal desalination are suspected to be less than those of other water supply alternatives, the uncertainty about potential site-specific impacts and their mitigation are large barriers to the application of coastal desalination in the United States. This uncertainty leads to stakeholder disagreements and a lengthy and costly planning and permitting process. For inland desalination, uncertainties remain about the sustainability of brackish groundwater resources and the environmental impacts from concentrate discharge to surface waters. Without rigorous scientific research to identify specific potential environmental impacts (or a lack of impacts), planners cannot assess the feasibility of desalination at a site or determine what additional mitigation steps are needed. Once potential impacts are clearly understood, research can be focused on developing approaches to minimize these impacts. The second goal focuses on the cost of desalination relative to the cost of other water supply alternatives. At present, costs are already low enough to make desalination an attractive option for some communities, especially where concentrate management costs are modest. In fact, desalination plants are being studied or implemented in at least 30 municipalities nationwide (GWI, 2007). The economic costs of desalination, however, as well as the costs of water supply alternatives, are locally variable. Costs are influenced by factors such as source water quality, siting considerations, potential environmental impacts, local regulations and permitting requirements, and available concentrate management op-
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Desalination: A National Perspective tions. Desalination remains a higher-cost alternative for many locations, and increasing awareness of potential environmental impacts is raising the costs of permitting and intake and outfall configurations in the United States. Inland communities considering brackish groundwater desalination may soon face more restrictions on surface water discharge and, therefore, will have fewer low-cost alternatives for concentrate management. Meanwhile, the future costs of energy are uncertain. If the total costs of desalination (including environmental costs) were reduced relative to other alternatives, desalination technology would become an attractive alternative to help address local water supply needs. STRATEGIC DESALINATION RESEARCH AGENDA The committee identified research topics as part of a strategic agenda to address the two long-term research goals articulated earlier. This agenda is driven by determination of what is necessary to make desalination a competitive option among other water supply alternatives. The agenda is broadly conceived, including research topics of clear interest to the public sector—and therefore of interest for federal funding—and research that might be most appropriately funded by private industry. The suggested research areas are described in detail below and are summarized in Box 8-1. Specific recommendations on the roles of federal and nonfederal organizations in funding the agenda are described in an upcoming section. BOX 8-1 Priority Research Areas The committee has identified priority research areas to help make desalination a competitive option among water supply alternatives for communities facing water shortages. These research areas, which are described in more detail in the body of the chapter, are summarized here. The highest priority topics are shown in bold. Some of this research may be most appropriately supported by the private sector. The research topics for which the federal government should have an interest—where the benefits are widespread and where no private-sector entities are willing to make the investments and assume the risk—are marked with asterisks. GOAL 1. Understand the environmental impacts of desalination and develop approaches to minimize these impacts relative to other water supply alternatives Assess environmental impacts of desalination intake and concentrate management approaches**
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Desalination: A National Perspective Conduct field studies to assess environmental impacts of brackish groundwater development** Develop protocols and conduct field studies to assess the impacts of concentrate management approaches in inland and coastal settings** Develop laboratory protocols for long-term toxicity testing of whole effluent to assess long-term impacts of concentrate on aquatic life** Assess the environmental fate and bioaccumulation potential of desalination-related contaminants** Develop improved intake methods at coastal facilities to minimize impingement of larger organisms and entrainment of smaller ones** Assess the quantity and distribution of brackish water resources nationwide** Analyze the human health impacts of boron, considering other sources of boron exposure, to expedite water-quality guidance for desalination process design** GOAL 2. Develop approaches to lower the financial costs of desalination so that it is an attractive option relative to other alternatives in locations where traditional sources of water are inadequate Improve pretreatment for membrane desalination Develop more robust, cost-effective pretreatment processes Reduce chemical requirements for pretreatment Improve membrane system performance Develop high-permeability, fouling-resistant, high-rejection, oxidant-resistant membranes Optimize membrane system design Develop lower-cost, corrosion-resistant materials of construction Develop ion-selective processes for brackish water Develop hybrid desalination processes to increase recovery Improve existing desalination approaches to reduce primary energy use Develop improved energy recovery technologies and techniques for desalination Research configurations and applications for desalination to utilize low-grade or waste heat** Understand the impact of energy pricing on desalination technology over time** Investigate approaches for integrating renewable energy with desalination** Develop novel approaches or processes to desalinate water in a way that reduces primary energy use** GOAL 1 and 2 Crosscuts Develop cost-effective approaches for concentrate management that minimize potential environmental impacts**
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Desalination: A National Perspective Research on Environmental Impacts The following research topics address Goal 1 to understand the environmental impacts of desalination and develop approaches to minimize those impacts relative to other water supply alternatives. Assess environmental impacts of desalination intake and concentrate management approaches As discussed in Chapter 5, the environmental impacts of desalination source water intake and concentrate management approaches are not well understood. Source water intakes for coastal desalination can create entrainment concerns with small organisms and impingement issues for larger organisms. For inland groundwater desalination, there are potential concerns regarding overpumping, water quality changes, and subsidence. The possible environmental impacts of concentrate management approaches range from effects on aquatic life in surface water discharges to the contamination of drinking water aquifers in poorly designed injection wells or ponds. Both site-specific studies and broad analyses of relative impacts would help communities weigh the alternatives for meeting water supply needs. The specific research needs are described as follows. 1a. Conduct field studies to assess environmental impacts of seawater intakes. Measurements and modeling of the extent of mortality of aquatic or marine organisms due to impingement and entrainment are needed. There have been numerous studies on such impacts of power plants, and extrapolation of such effects to desalination facilities should be performed. 1b. Conduct field studies to assess environmental impacts of brackish groundwater development. The general environmental interactions between wetlands, freshwater, and brackish aquifers for inland sources have not been documented under likely brackish water development scenarios. While site-specific evaluation of any location will be necessary for developing a brackish water resource, the lack of synthesized information is an impediment to the use of this resource for smaller communities with limited resources. 1c. Develop protocols and conduct field studies to assess the impacts of concentrate management approaches in inland and coastal settings. Comprehensive studies analyzing impacts of concentrate discharge at marine, estuarine, and inland desalination locations are needed.
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Desalination: A National Perspective Adequate site-specific baseline studies on potential biological and ecological effects are necessary prior to the development of desalination facilities because biological communities in different geographic areas will have differential sensitivity, but a comprehensive synthesis would be valuable once several in-depth studies have been conducted. Protocols should be developed to define the baseline and operational monitoring, reference sites, lengths of transects, and sampling frequencies. Planners would benefit from clear guidance on appropriate monitoring and assessment protocols. Environmental data should be collected for at least 1 year in the area of the proposed facility before a desalination plant with surface water concentrate discharge comes online so that sufficient baseline data on the ecosystem are available with which to compare postoperating conditions. Once a plant is in operation, monitoring of the ecological communities (especially the benthic community) receiving the concentrate should be performed periodically for at least 2 years at multiple distances from the outflow pipe and compared to reference sites. For inland settings, additional regional hydrogeology research is needed on the distribution, thickness, and hydraulic properties of formations that could be used for disposal of concentrate via deep-well injection. Much information is already available about the potential for deep-well injection in states such as Florida and Texas, although suitable geologic conditions may exist in other states as well. Inventories of industrial and commercial brine-disposal wells and producing and abandoned oil fields should be synthesized and used to develop a suitable protocol for further hydrogeological investigations, as appropriate. This research would provide valuable assistance to small communities that typically do not have the resources available to support extensive hydrogeological investigations. 1d. Develop laboratory protocols for long-term toxicity testing of whole effluent to assess long-term impacts of concentrate on aquatic life. Standard acute toxicity tests as defined by the U.S. Environmental Protection Agency (EPA) are generally 96 hours in duration and use larval or juvenile stages of certain fish and invertebrate species with a series of effluent dilutions and a control. The end point is whether the test organisms survive or not. Chronic tests, according to EPA, are typically 7 days in duration when using larval stages of fish and invertebrate species, and the end points of the tests are sublethal, such as growth reduction. Typical chronic toxicity protocols were designed for testing municipal or industrial wastewater treatment plant effluent, which typically contains higher levels of toxic chemicals than the concentrate from desalination plants. To assess the impacts of desalination effluent, a protocol should be developed to analyze the longer-term effects (over whole life cycles) on organisms that live in the vicinity of desalination plants (as opposed
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Desalination: A National Perspective to the standard species used in EPA-required toxicity testing). These laboratory-based tests should then be used to examine the impacts of whole effluent (and various dilutions) from different desalination plants on a variety of different taxa at numerous representative sites from key ecological regions. 1e. Assess the environmental fate and bioaccumulation potential of desalination-related contaminants. Desalination concentrate contains more than just salts and may include various chemicals that are used in pretreatment and membrane cleaning, antiscaling and antifoulant additives, and metals that may leach from corrosion. Some of these chemicals (e.g., antifoulants, copper leached from older thermal desalination plants) or chemical by-products (e.g., trihalomethanes produced as a result of pretreatment with chlorine) are likely to bioaccumulate in organisms. Investigations into the loading and environmental fate of desalination-related chemicals should be included in modeling and monitoring programs. The degree to which various chemicals biodegrade or accumulate in sediments should also be investigated. High priority should be given to polymer antiscalants, such as polycarbonic acids and polyphosphate, which may increase primary productivity. Corrosion-related metals and disinfection by-products should also be investigated. In conjunction with the field studies described earlier, representative species, preferably benthic infauna along the transects and from the reference (control) site, should be analyzed for bioaccumulative contaminants. Because little is known about the potential of some other desalination chemicals that can be discharged in concentrate to bioaccumulate (e.g., polyphosphate, polycarbonic acid, polyacrylic acid, polymaleic acid), research should be conducted into their toxicity and bioaccumulation potential. Develop improved intake methods at coastal facilities to minimize impingement of larger organisms and entrainment of smaller ones Although intake and screen technology is rapidly developing, continued research and development is needed in the area of seawater intakes to develop cost-effective approaches that minimize the impacts of impingement and entrainment for coastal desalination facilities. Current technology development has focused on subsurface intakes and advanced screens or curtains, and these recent developments should be assessed to determine the costs and benefits of the various approaches. Other innovative concepts could also be considered that might deter marine life from entering intakes.
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Desalination: A National Perspective Assess the quantity and distribution of brackish water resources nationwide Sustainable development of inland brackish water resources requires maps and synthesized information on total dissolved solids of the groundwater, types of dominant solutes (e.g., NaCl, CaSO4), thickness, and depth to brackish water. The only national map of brackish water resources available (Feth, 1965; Figure 1-1) simply shows depth to saline water. Newer and better solute chemistry data collected over the past 40 years exist in the files of private, state, and federal offices but are not generally organized for use in brackish water resources investigations. Using the aforementioned information, basin analyses, analogous to the U.S. Geological Survey Regional Aquifer System Analysis program for freshwater, could be developed, emphasizing regions facing near-term water scarcity concerns. These brackish water resource investigations could also be conducted at the state level. The data, once synthesized, could be utilized for desalination planning as well as for other water resources and commercial development scenarios. Analyze the human health impacts of boron, considering other sources of boron exposure, to expedite water-quality guidance for desalination process design Typical single-pass reverse osmosis (RO) desalination processes do not remove all the boron in seawater; thus, boron can be found at milligram-per-liter levels in the finished water. Boron can be controlled through treatment optimization, but that treatment has an impact on the cost of desalination. A range of water quality levels (0.5 to 1.4 mg/L) have been proposed as protective of public health based on different assumptions in the calculations. Because of the low occurrence of boron in most groundwater and surface water, the EPA has decided not to develop a maximum contaminant level for boron and has encouraged affected states to issue guidance or regulations as appropriate (see Chapter 5). Additional analysis of existing boron toxicity data is needed, considering other possible sources of boron exposure in the United States, to support guidance for desalination process design that will be suitably protective of human health.
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Desalination: A National Perspective Research to Lower the Costs of Desalination The following research topics address Goal 2 to develop approaches to lower the costs of desalination so that it is an attractive option relative to other alternatives in locations where traditional sources of water are inadequate. As a broadly conceived agenda, some of this research may be most appropriately supported by the private sector. The appropriate roles of governmental and nongovernmental entities to fund the research agenda are discussed later in the chapter. Improve pretreatment for membrane desalination Pretreatment is necessary to remove potential foulants from the source water, thereby ensuring sustainable operation of the RO membranes at high product water flux and salt rejection. Research to improve the pretreatment process is needed that would develop alternative, cost-effective approaches. 5a. Develop more robust, cost-effective pretreatment processes. Membrane fouling is one of the most problematic issues facing seawater desalination. Forms of fouling common with RO membranes are organic fouling, scaling, colloidal fouling, and biofouling. All forms of fouling are caused by interactions between the foulant and the membrane surface. Improved pretreatment that minimizes these interactions will reduce irreversible membrane fouling. Alteration of solution characteristics can improve the solubility of the foulants, preventing their precipitation or interaction with the membrane surface. Such alteration could be chemical, electrochemical, or physical in nature. Membranes such as microfiltration (MF) and ultrafiltration (UF) have several advantages over traditional pretreatment (e.g., conventional sand filtration) because they have a smaller footprint, are more efficient in removing smaller foulants, and provide a more stable influent to the RO membranes. Additional potential benefits of MF or UF pretreatment are increased flux, increased recovery, longer membrane life, and decreased cleaning frequency. More research is necessary in order to optimize the pretreatment membranes for more effective removal of foulants to the RO system, to reduce the fouling of the pretreatment membranes, and to improve configuration of the pretreatment membranes to maximize cost reduction. 5b. Reduce chemical requirements for pretreatment. Antiscalants, coagulants, and oxidants (such as chlorine) are common chemicals
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Desalination: A National Perspective applied in the pretreatment steps for RO membranes. Although these chemicals are added to reduce fouling, they add to the operational costs, can reduce the operating life of membranes, and have to be disposed of properly or they can adversely impact aquatic life (see Chapter 5). Antiscalants may also enhance biofouling, so alternative formulations or approaches should be examined. Research is needed on alternative formulations or approaches (including membrane pretreatment) to reduce the chemical requirements of the pretreatment process, both to reduce overall cost and to decrease the environmental impacts of desalination. Improve membrane system performance Sustainable operation of the RO membranes at the designed product water flux and salt rejection is a key to the reduction of desalination process costs. In addition to effective pretreatment, research to optimize the sustained performance of the RO membrane system is needed. 6a. Develop high-permeability, fouling-resistant, high-rejection, oxidant-resistant membranes. New membrane designs could reduce the treatment costs of desalination by improving membrane permeability and salt rejection while increasing resistance to fouling and membrane oxidation. Current membrane research to reduce fouling includes altering the surface charge, increasing hydrophilicity, adding polymers as a barrier to fouling, and decreasing surface roughness. Oxidant-resistant membranes enable feedwater to maintain an oxidant residual that will reduce membrane fouling due to biological growth. Current state-of-the-art thin-film composite desalination membranes are polyamide based and therefore are vulnerable to damage by chlorine or other oxidants. Thus, when an oxidant such as chlorine is added to reduce biofouling, dechlorination is necessary to prevent structural damage. Additionally, trace concentrations of chlorine may be present in some feedwaters. Cellulose-derivative RO membranes have much higher chlorine tolerance; however, these membranes have a much lower permeability than thin-film composite membranes and operate under a narrower pH range. Therefore, there is a need to increase the oxidant tolerance of the higher-permeability membranes. Lower risk of premature membrane replacement equates to overall lower operating costs. Past efforts to synthesize RO membranes with high permeability often resulted in reduced rejection and selectivity. There is a need to develop RO membranes with high permeability without sacrificing selectivity or rejection efficiency. Recent research on utilizing nanomaterials, such as carbon nanotubes, as a separation barrier suggest the possibility
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Desalination: A National Perspective of obtaining water fluxes much higher than that of traditional polymeric membranes. The development of membranes that are more resistant to degradation from exposure to cleaning chemicals will extend the useful life of a membrane module. The ability to clean membranes more frequently can also decrease energy usage because membrane fouling results in higher differential pressure loss through the modules. By extending the life of membrane modules, the operating and maintenance cost will be reduced by the associated reduction in membrane replacements required. 6b. Optimize membrane system design. With the development of high-flux membranes and larger-diameter membrane modules, new approaches for optimal RO system design are needed to avoid operation under thermodynamic restriction (see Chapter 4) and to ensure equal distribution of flux between the leading and tail elements of the RO system. The key variables for the system design will involve the choice of optimal pressure, the number of stages, and number and size of membrane elements at each stage. An optimal system configuration may also involve hybrid designs where one type of membrane (e.g., intermediate flux, highly fouling-resistant) is used in the leading elements followed by high-flux membranes in the subsequent elements. Fouling can be mitigated by maintaining high crossflow velocity; thus, fouling-resistant membranes may be better served in the downstream positions where lower crossflow velocity is incurred. Thus, additional engineering research on membrane system design is needed to optimize performance with the objective of reducing costs. 6c. Develop lower-cost, corrosion-resistant materials of construction. The duration of equipment life in a desalination plant directly relates to the total costs of the project. Saline and brackish water plants are considered to be a corrosive environment due to the high levels of salts in the raw water. The development and utilization of corrosion-resistant materials will minimize the frequency of equipment or appurtenance replacement, which can significantly reduce the total project costs. 6d. Develop ion-selective processes for brackish water. Some slightly brackish waters could be made potable simply though specific removal of certain contaminants, such as nitrate or arsenite, while removing other ions such as sodium, chloride, and bicarbonate at a lower rate. High removal rates of all salts are not necessary for such waters. Ion-specific separation processes, such as an ion-selective membrane or a selective ion-exchange resin, should be able to produce potable water at much lower energy costs than those processes that fully desalinate the
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Desalination: A National Perspective source water. Ion-selective removal would also create fewer waste materials requiring disposal. Ion-selective processes would be useful for mildly brackish groundwater sources with high levels of nitrate, uranium, radium, or arsenic. Such an ion-selective process could also be used to optimize boron removal following RO desalination of seawater. 6e. Develop hybrid desalination processes to increase recovery. Overall product water recovery in a desalination plant can be increased through the serial application of more than one desalination process. For example, an RO process could be preceded by a “tight” nanofiltration process, allowing the RO to operate at a higher recovery than it could with less aggressive pretreatment. Other options could be devised, including hybrid thermal and membrane processes to increase the overall recovery of the process. As noted in Chapter 4, the possible hybrid combinations of desalination processes are limited only by ingenuity and identification of economically viable applications. Hybridization also offers opportunities for reducing desalination production costs and expanding the flexibility of operations, especially when co-located with power plants, but hybridization also increases plant complexity and raises challenges in operation and automation. Improve existing desalination approaches to reduce primary energy use Energy is one of the largest annual costs in the desalination process. Thus, research to improve the energy efficiency of desalination technologies could make a significant contribution to reducing costs. 7a. Develop improved energy recovery technologies and techniques for desalination. Membrane desalination is an energy-intensive process compared to treatment of freshwater sources. Modern energy recovery devices operate at up to 96 percent energy recovery (see Chapter 4), although these efficiencies are lower at average operating conditions. The energy recovery method in most common use today is the energy recovery (or Pelton) turbine, which achieves about 87 percent efficiency. Many modern plants still use Pelton wheels because of the higher capital cost of isobaric devices. Thus, opportunities exist to improve recovery of energy from the desalination concentrate over a wide operating range and reduce overall energy costs. 7b. Research configurations and applications for desalination to utilize low-grade or waste heat. Industrial processes that produce waste or low-grade heat may offer opportunities to lower the operating cost of
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Desalination: A National Perspective the desalination process if these heat sources are co-located with desalination facilities (see Box 4-8). Low-grade heat can be used as an energy source for desalination via commercially available thermal desalination processes. Hybrid membrane-thermal desalination approaches offer additional operational flexibility and opportunities for water-production cost savings. Research is needed to examine configurations and applications of current technologies to utilize low-grade or waste heat for desalination. 7c. Understand the impact of energy pricing on existing desalination technology over time. Energy is one of the largest components of cost for desalination, and future changes in energy pricing could significantly affect the affordability of desalination. Research is needed to examine to what extent the economic and financial feasibility of desalination may be threatened by the uncertain prospect of energy price increases in the future for typical desalination plants in the United States. This research should also examine the costs and benefits of capital investments in renewable energy sources. 7d. Investigate approaches for integrating renewable energy with desalination. Renewable energy sources could help mitigate future increases in energy costs by providing a means to stabilize energy costs for desalination facilities while also reducing the environmental impacts of water production. Research is needed to optimize the potential for coupling various renewable energy applications with desalination. Develop novel approaches or processes to desalinate water in a way that reduces primary energy use Because the energy of RO is only twice the minimum energy of desalination, even novel technologies are unlikely to create step change (>25 percent) reductions in absolute energy consumption compared to the best current technology (see, e.g., Appendix A). Instead, substantial reductions in the energy costs of desalination are more likely to come through the development of novel approaches or processes that optimize the use of low-grade heat. Several innovative desalination technologies that are the focus of ongoing research, such as forward osmosis, dewvaporation, and membrane distillation, have the capacity to use low-grade heat as an energy source. Research into the specific incorporation of waste or low-grade heat into these or other innovative processes could greatly reduce the amount of primary energy required for desalination and, thus, overall desalination costs.
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Desalination: A National Perspective Crosscutting Research Research topics in this category benefit both Goal 1, for environmental impacts, and Goal 2, for lowering the cost of desalination. Develop cost-effective approaches for concentrate management that minimize potential environmental impacts Research objectives related to concentrate management are crosscutting, because they address both the need to understand and minimize environmental impacts and the need to reduce the total cost of desalination. For coastal concentrate management, research is needed to develop improved diffuser technologies and subsurface injection approaches and to examine their costs and benefits relative to current disposal alternatives. The high cost of inland concentrate management inhibits inland brackish water desalination. Low- to moderate-cost concentrate management alternatives (i.e., subsurface injection, land application, sewer discharge, and surface water discharge) can be limited by the salinity of the concentrate and by location and climate factors; in some scenarios all of these options may be restricted by site-specific conditions, leaving zero liquid discharge (ZLD) as the only alternative for consideration. ZLD options currently include evaporation ponds and energy-intensive processes, such as brine concentrators or crystallizers, followed by landfilling. These options have high capital or operating costs. Research to improve recovery in the desalination process and thereby minimize the initial volume of concentrate could enhance the practical viability of several concentrate management options for inland desalination. This is particularly true for the concentrate management options that are characterized by high costs per unit volume of the concentrate flow treated and for approaches that are not applicable to large concentrate flows, such as thermal evaporation or evaporation ponds. Advancements are also needed that reduce the capital costs and improve the energy efficiency of thermal evaporation processes. Conventional concentrate management options that involve simple equipment are not likely to see significant cost reductions through additional research. The reuse of high-salinity concentrates and minerals extracted from them should be further explored and developed to help mitigate environmental impacts while generating revenues that can help offset concentrate management costs. Possibilities include selective precipitation of marketable salts, irrigation of salt-tolerant crops, supplements for animal dietary needs, dust suppressants, stabilizers for road base construction, or manufacture of lightweight fire-proof building materials. Studies are
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Desalination: A National Perspective necessary to determine the most feasible uses and to develop ways to prepare the appropriate product for various types of reuse. For all possible uses, site-specific limitations and local and state regulations will need to be considered. Because the transportation costs greatly affect the economics of reuse, a market analysis would also be needed to identify areas in the United States that could reasonably utilize products from desalination concentrate. Highest Priority Research Topics All of the topics identified are considered important, although three topics (1, 2, and 9 above) were deemed to be the highest priority research topics: (1) assessing the environmental impacts of desalination intake and concentrate management approaches, (2) developing improved intake methods to minimize impingement and entrainment, and (3) developing cost-effective approaches for concentrate management that minimize environmental impacts. These three research areas are considered the highest priorities because this research can help address the largest barriers (or showstoppers) to more widespread use of desalination in the United States. Uncertainties about potential environmental impacts will need to be resolved and cost-effective mitigation approaches developed if desalination is to be more widely accepted. Research to develop cost-effective approaches for concentrate management is critical to enable more widespread use of desalination technologies for inland communities. As noted in Chapter 4, the cost of concentrate management can double or triple the cost of the desalination for some inland communities. Research may also reduce the costs of desalination. Any cost improvement will help make desalination an attractive option for communities addressing water shortages. However, the committee does not view these process cost issues as the major limitation to the application of desalination in the United States today. IMPLEMENTING THE RESEARCH AGENDA In the previous section, the committee proposed a broad research agenda that, if implemented, should improve the capacity of desalination to meet future water needs in the United States by further examining and addressing its environmental impacts and reducing its costs relative to other water supply alternatives. Implementing this agenda requires federal leadership, but its success depends on participation from a range of entities, including federal, state, and local governments, nonprofit or-
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Desalination: A National Perspective ganizations, and the private sector. A strategy for implementing the research agenda is suggested in the following section. This section also includes suggestions for funding the agenda and the appropriate roles of government and nongovernmental entities. Supporting the Desalination Research Agenda A federal role is appropriate for research that provides a “public good.” Specifically, the federal government should have an interest in funding research where the benefits are widespread but where no private-sector entities are willing to make the investment and assume the risks. Thus, for example, research that results in significant environmental benefits should be in the federal interest because these benefits are shared by the public at large and cannot be fully captured by any entrepreneur. Federal investment is also important where it has “national significance”—where the issues are of large-scale concern; they are more than locally, state-, or regionally specific; and the benefits accrue to a large swath of the public. Based on the aforementioned criteria, the proposed research agenda contains many topic items that should be in the federal interest (see topics marked with asterisks in Box 8-1). The research topics in support of Goal 1 (see Box 8-1) are directed at environmental issues that are largely “public good” issues. Some of the needed environmental research will, by nature, be site-specific, and purely site-specific research is not of great federal interest. Thus, there is a clear role for state and local agencies to support site-specific research. The federal government, however, should have an interest in partnering with local communities to conduct more extensive field research from which broader conclusions of environmental impacts can be drawn or which would significantly contribute to a broader meta-analysis. This meta-analysis could especially benefit small water supply systems. Also, there should be federal interest in establishing general protocols for field evaluations and chronic bioassays that could then be adapted for site-specific studies. The research needed to support the attainment of Goal 2 includes several topics that are clearly in the federal interest, as defined earlier. These include efforts to reduce prime energy use, to integrate renewable energy resources within the total energy picture and increase reliance upon them, and to understand the impacts of energy pricing on the future of desalination (see highlighted topics in Box 8-1). However, Goal 2 also includes a number of research topics that may be more appropriately funded by the private sector or nongovernmental organizations, assuming that these entities are willing to assume the risks of the research investment. Indeed, private industry already spends far more on research and
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Desalination: A National Perspective development for desalination than the federal government (see Chapter 2) and is already making substantial progress in the improvement of existing membrane performance, developing better pretreatment alternatives, and developing improved energy recovery devices. To avoid duplication and to optimize available research funding, government programs should focus instead on research and development with widespread possible benefits that would otherwise go unfunded because private industry is unwilling to make the investment. Finally, the crosscutting topic to develop cost-effective methods of managing concentrates for inland communities, which impacts Goals 1 and 2, is also in the federal interest. Federal Research Funding The optimal level of federal investment in desalination research is inherently a question of public policy. Although the decision should be informed by science, it is not—at its heart—a scientific decision. However, several conclusions emerged from the committee’s analysis of current research and development funding (see Chapter 2) that suggest the importance of strategic integration of the research program. The committee concluded that there is no integrated and strategic direction to the federal desalination research and development efforts. Continuation of a federal program of research dominated by congressional earmarks and beset by competition between funding for research and funding for construction will not serve the nation well and will require the expenditure of more funds than necessary to achieve specified goals. To ensure that future federal investments in desalination research are integrated and prioritized so as to address the two major goals identified in this report, the federal government will need to develop a coordinated strategic plan that utilizes the recommendations of this report as a basis. It is beyond the committee’s scope to recommend specific plans for improving coordination among the many federal agencies that support desalination research. Instead, responsibility for developing the plan should rest with the Office of Science and Technology Policy’s (OSTP’s) National Science and Technology Council (NSTC) because “this Cabinet-level Council is the principal means within the executive branch to coordinate science and technology policy across the diverse entities that make up the Federal research and development enterprise.”1 For example, the NSTC’s Subcommittee on Water Availability and Quality has member-ship representing more than 20 federal agencies and recently released “A Strategy for Federal Science and Technology to Support Water Avail- 1 For more information, see http://www.ostp.gov/nstc/index.html.
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Desalination: A National Perspective ability and Quality in the United States” (SWAQ, 2007). Representatives of the National Science Foundation, the Bureau of Reclamation, the Environmental Protection Agency, the National Oceanographic and Atmospheric Administration, the Office of Naval Research, and the Department of Energy should participate fully in the development of the strategic federal plan for desalination research and development. Five years into the implementation of this plan, the OSTP should evaluate the status of the plan, whether goals have been met, and the need for further funding. A coordinated strategic plan governing desalination research at the federal level along with effective implementation of the research plan will be the major determinants of federal research productivity in this endeavor. The committee cannot emphasize strongly enough the importance of a well-organized, well-articulated strategically directed effort. In the absence of any or all of these preconditions, federal investment will yield less than it could. Therefore, a well-developed and clearly articulated strategic research plan, as called for above, should be a precondition for any new federal appropriations. Initial federal appropriations on the order of recent spending on desalination research (total appropriations of about $25 million annually, as in fiscal years 2005 and 2006) should be sufficient to make good progress toward the overall research goals if the funding is strategically directed toward the proposed research topics as recommended in this report. Annual federal appropriations of $25 million, properly allocated, should be sufficient to have an impact in the identified priority research areas, given the context of expected state and private-sector funding. This level of federal funding is also consistent with NRC (2004a), which recommended annual appropriations of $700 million for research supporting the nation’s entire water resources research agenda. Reallocation of current spending will be necessary to address topics that are currently underfunded. If current research funding is not reallocated, the overall desalination research and development budget will need to be enhanced. Nevertheless, support for the research agenda stated here should not come at the expense of other high-priority water resource research topics, such as those identified in Confronting the Nation’s Water Problems: The Role of Research (NRC, 2004a). Environmental research should be emphasized up front in the research agenda. At least 50 percent of the federal funding for desalination research should initially be directed toward environmental research. Environmental research, including Goal 1 and the Goal 1 and 2 crosscuts, should be addressed, because these have the potential for the greatest impact in overcoming current roadblocks for desalination and making desalination an attractive water supply alternative. Research funding in support of Goal 2 should be directed strategically toward research topics that are likely to make improvements against benchmarks set by the best
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Desalination: A National Perspective current technologies for desalination. The best available technologies for desalination at the time of this writing are benchmarked in Chapter 4. Research proposals should make the case as to how and to what degree the proposed research can advance the state of the art in desalination. An emphasis should be placed on energy benchmarks because reductions in energy result in overall cost savings and have environmental benefits. The majority of the federal funding directed toward Goal 2 should support projects that are in the public interest and would not otherwise be privately funded (see Box 8-1), such as some high-risk and long-term research initiatives (e.g., developing novel desalination processes that sharply reduce the primary energy use). Although private industry does make modest investments in high-risk research, it is frequently reluctant to invest in research in the earliest stage of technology creation, when there is extremely low likelihood of success even though there are large potential benefits. The effectiveness with which federal funds are spent will also depend on certain critical implementation steps, which are outlined in the following section. Proposal Announcement and Selection Based on available funding, the opportunity to announce requests for proposals exists for federal agencies, such as the Bureau of Reclamation or the National Science Foundation, or other research institutions that explicitly target one or more research objectives. The principal funding agency should announce a request for proposals as widely as possible to scientists and engineers in municipal and federal government, academia, and private industry. At present, the desalination community is relatively small, but collectively there is a great deal of expertise across the world. International desalination experts and others from related areas of research should be encouraged and given the opportunity to offer innovative research ideas that have the potential to significantly advance the field. Thus, the request for proposals should extend to federal agencies, national laboratories, other research institutions, utilities, and the private sector. Since innovation cannot be preassigned, broad solicitations for proposals should include a provision for unsolicited investigator-initiated research proposals. To achieve the objectives of the research agenda, proposals should be selected through a rigorous independent peer-review process (NRC, 2002b) irrespective of the agency issuing the request for proposals. A rotating panel of independent, qualified reviewers should be appointed based on their relevant expertise in the focal areas. The process should
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Desalination: A National Perspective allow for the consideration and review of unsolicited proposals, as long as their research goals meet the overall research goals. Proposal funding should be based on the quality of the proposed work, the degree to which the proposed research can advance the state of the art in desalination or otherwise contribute toward the research goals, prior evidence of successful research, and the potential for effective publication or dissemination of the research findings. CONCLUSIONS AND RECOMMENDATIONS A strategic national research agenda has been conceived that centers around two overarching strategic goals for further research in desalination: (1) to understand the environmental impacts of desalination and develop approaches to minimize these impacts relative to other water supply alternatives and (2) to develop approaches to lower the financial costs of desalination so that it is an attractive option relative to other alternatives in locations where traditional sources of water are inadequate. A research agenda is proposed in this chapter in support of these two goals (see Box 8-1). Several recommendations for implementing the proposed research agenda follow. A coordinated strategic plan should be developed to ensure that future federal investments in desalination research are integrated and prioritized and address the two major goals identified in this report. The strategic application of federal funding for desalination research can advance the implementation of desalination technologies in areas where traditional sources of water are inadequate. Responsibility for developing the plan should rest with the OSTP, which should use the recommendations of this report as a basis for plan development. Initial federal appropriations on the order of recent spending on desalination research (total appropriations of about $25 million annually) should be sufficient to make good progress toward these goals, when complemented by ongoing nonfederal and private-sector desalination research, if the funding is directed toward the proposed research topics as recommended in this chapter. Reallocation of current federal spending will be necessary to address currently underfunded topics. If current federal research and development funding is not reallocated, new appropriations will be necessary. However, support for the research agenda stated here should not come at the expense of other high-priority water resource research topics. Five years into the implementation of this plan, the OSTP should evaluate the status of the plan, whether goals have been met, and the need for further funding.
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Desalination: A National Perspective Environmental research should be emphasized up front when implementing the research agenda. Uncertainties regarding environmental impacts and ways to mitigate these impacts are one of the largest hurdles to implementation of desalination in the United States, and research in these areas has the greatest potential for enabling desalination to help meet future water needs in communities facing water shortages. This environmental research includes work to understand environmental impacts of desalination intakes and concentrate management, the development of improved intake methods to minimize impingement and entrainment, and cost-effective concentrate management technologies. Research funding in support of reducing the costs of desalination (Goal 2) should be directed strategically toward research topics that are likely to make improvements against benchmarks set by the best current technologies for desalination. Because the private sector is already making impressive strides toward Goal 2, federal research funding should emphasize the long-term and high-risk research that may not be attempted by the private sector and that is in the public interest, such as research on novel technologies that significantly reduce prime energy use. Wide dissemination of requests for proposals to meet the goals of the research agenda will benefit the quality of research achieved. Requests for proposals should extend to federal agencies, national laboratories, research institutions, utilities, other countries, and the private sector. Investigator-driven research through unsolicited proposals should be permitted throughout the proposal process. Proposals should be peer-reviewed and based on quality of research proposed, the potential contribution, prior evidence of successful research, and effective dissemination. |
Table of Contents
PC-style floppy disks work mostly like other disk devices like hard disks, except that you need to low-level format them first. To use an common 1440 KB floppy in the first floppy drive, first (as root) format it:
#fdformat -f /dev/rfd0a
Then create a single partition on the disk using disklabel(8):
#disklabel -rw /dev/rfd0a floppy3
Creating a small filesystem optimized for space:
#newfs -m 0 -o space -i 16384 -c 80 /dev/rfd0a
Now the floppy disk can be mounted like any other disk. Or if you already have a floppy disk with an MS-DOS filesystem on it that you just want to access from NetBSD, you can just do something like this:
#mount -t msdos /dev/fd0a /mnt
However, rather than using floppies like normal (bigger) disks, it is often more convenient to bypass the filesystem altogether and just splat an archive of files directly to the raw device. E.g.:
#tar cvfz /dev/rfd0a file1 file2 ...
A variation of this can also be done with MS-DOS floppies using
sysutils/mtools package which
has the benefit of not going through the kernel buffer cache
and thus not being exposed to the danger of the floppy being
removed while a filesystem is mounted on it.
See if your system has a ZIP drive:
dmesg | grep -i zipsd0 at atapibus0 drive 1: <IOMEGA ZIP 100 ATAPI, , 14.A> type 0 direct removable
Seems it has one, and it's recognized as sd0, just like any SCSI disk. The fact that the ZIP here is an ATAPI one doesn't matter - a SCSI ZIP will show up here, too. The ZIP is marked as "removable", which means you can eject it with:
Insert ZIP disk
Check out what partitions are on the ZIP:
#/dev/rsd0d: type: ATAPI ... 8 partitions:
#size offset fstype [fsize bsize cpg] d: 196608 0 unused 0 0 # (Cyl. 0 - 95) h: 196576 32 MSDOS # (Cyl. 0*- 95) disklabel: boot block size 0 disklabel: super block size 0
is the whole disk, as usual on i386.
is what you want, and you can see it's a msdos filesystem even.
Hence, use /dev/sd0h to access the zip's partition.
mount -t msdos /dev/sd0h /mnt
Access your files:
ls -la /mnttotal 40809 drwxr-xr-x 1 root wheel 16384 Dec 31 1979 . drwxr-xr-x 28 root wheel 1024 Aug 2 22:06 .. -rwxr-xr-x 1 root wheel 1474560 Feb 23 1999 boot1.fs -rwxr-xr-x 1 root wheel 1474560 Feb 23 1999 boot2.fs -rwxr-xr-x 1 root wheel 548864 Feb 23 1999 boot3.fs -rwxr-xr-x 1 root wheel 38271173 Feb 23 1999 netbsd19990223.tar.gz
Unmount the ZIP:
Eject the ZIP:
Data CDs can contain anything from programs, sound files (MP3, wav), movies (MP3, QuickTime) to source code, text files, etc. Before accessing these files, a CD must be mounted on a directory, much like hard disks are. Just as hard disks can use different filesystems (ffs, lfs, ext2fs, ...), CDs have their own filesystem, "cd9660". The NetBSD cd9660 filesystem can handle filesystems without and with Rockridge and Joliet extensions.
CD devices are named /dev/cd0a for both SCSI and IDE (ATAPI).
With this information, we can start:
See if your system has some CD drive:
dmesg | grep ^cdcd0 at atapibus0 drive 0: <CD-R/RW RW8040A, , 1.12> type 5 cdrom removable cd0: 32-bit data port cd0: drive supports PIO mode 4, DMA mode 0 cd0(pciide0:1:0): using PIO mode 0, DMA mode 0 (using DMA data transfers)
We have one drive here, "cd0". It is an IDE/ATAPI drive, as it is found on atapibus0. Of course the drive (rather, its medium) is removable, i.e., you can eject it. See below.
Insert a CD
Mount the CD manually:
mount -t cd9660 /dev/cd0a /mnt
This command shouldn't print anything. It instructs the system to mount the CD found on /dev/cd0a on /mnt, using the "cd9660" filesystem. The mountpoint "/mnt" must be an existing directory.
Check the contents of the CD:
ls /mntINSTALL.html INSTALL.ps TRANS.TBL boot.catalog INSTALL.more INSTALL.txt binary installation
Everything looks fine! This is a NetBSD CD, of course. :)
Unmount the CD:
If the CD is still accessed (e.g. some other shell's still "cd"'d into it), this will not work. If you shut down the system, the CD will be unmounted automatically for you, there's nothing to worry about there.
Making an entry in /etc/fstab:
If you don't want to type the full "mount" command each time, you can put most of the values into a line in /etc/fstab:
# Device mountpoint filesystem mount options /dev/cd0a /cdrom cd9660 ro,noauto
Make sure that the mountpoint,
/cdrom in our
Now you can mount the cd with the following command:
Access and unmount as before.
The CD is not mounted at boot time due to the "noauto" mount option - this is useful as you'll probably not have a CD in the drive all the time. See mount(8) and mount_cd9660(8) for some other useful options.
Eject the CD:
If the CD is still mounted, it will be unmounted if possible, before being ejected.
Use mscdlabel(8) to add all sessions to the CDs
then use the appropriate device node to mount the session you want.
You might have to create the corresponding device nodes in
mscdlabel cd1track (ctl=4) at sector 142312 adding as 'a' track (ctl=4) at sector 0 adding as 'b'
ls -l /dev/cd1bls: /dev/cd1b: No such file or directory
ls -l cd1*brw-r----- 1 root operator 6, 8 Mar 18 21:55 cd1a brw-r----- 1 root operator 6, 11 Mar 18 21:55 cd1d
mknod cd1b b 6 9
Make sure you fix the permissions of any new device
nodes you create:
ls -l cd1*brw-r----- 1 root operator 6, 8 Mar 18 21:55 cd1a brw-r--r-- 1 root wheel 6, 9 Mar 18 22:23 cd1b brw-r----- 1 root operator 6, 11 Mar 18 21:55 cd1d
chgrp operator cd1b
chmod 640 cd1b
ls -l cd1*brw-r----- 1 root operator 6, 8 Mar 18 21:55 cd1a brw-r----- 1 root operator 6, 9 Mar 18 22:24 cd1b brw-r----- 1 root operator 6, 11 Mar 18 21:55 cd1d
Now you should be able to mount it.
mount /dev/cd1b /mnt
By default, NetBSD only allows "root" to mount a filesystem. If you want any user to be able to do this, perform the following steps:
Give groups and other the access rights to the device.
#chmod go+rw /dev/cd0a
Ask NetBSD to let users mounting filesystems.
#sysctl -w vfs.generic.usermount=1
Note that this works for any filesystem and device, not only for CDs with a ISO 9660 filesystem.
To perform the mount operation after these commands, the user must own the mount point. So, for example:
$mount -t cd9660 -o nodev,nosuid /dev/cd0a `pwd`/cdrom
The mount options
nosuid are mandatory from NetBSD 4.0 on. They are
not necessary on NetBSD 3.x systems.
Sometimes, it is interesting to mount an ISO9660 image file before you burn the CD; this way, you can examine its contents or even copy files to the outside. If you are a Linux user, you should know that this is done with the special loop filesystem. NetBSD does it another way, using the vnode pseudo-disk.
We will illustrate how to do this with an example. Suppose you have an ISO image in your home directory, called "mycd.iso":
Start by setting up a new vnode, "pointing" to the ISO file:
vnconfig -c vnd0 ~/mycd.iso
Now, mount the vnode:
mount -t cd9660 /dev/vnd0a /mnt
Yeah, image contents appear under
Go to that
directory and explore the image.
When you are happy, you have to umount the image:
And at last, deconfigure the vnode:
vnconfig -u vnd0
Note that these steps can also be used for any kind of file that contains a filesystem, not just ISO images.
To play MPEG Video streams as many DVD players can play them
under NetBSD, mount the CD as you would do with any normal (data)
CD (see Section 13.3, “Reading data CDs with NetBSD”), then use the
package to play the mpeg files stored on the CD.
There are two ways to handle audio CDs:
Tell the CD drive to play to the headphone or to a
soundcard, to which CDROMs are usually connected
internally. Use programs like cdplay(1),
multimedia/kdemultimedia3 package, mixer
programs like mixerctl(1),
the Curses based
or kmix, which is part of
This usually works well on both SCSI and IDE (ATAPI) CDROMs, CDRW and DVD drives.
To read ("rip") audio tracks in binary form without going through digital->analog conversion and back. There are several programs available to do this:
For most ATAPI, SCSI and several proprietary
CDROM drives, the
audio/cdparanoia package can be
used. With cdparanoia the data can be saved to a file or
directed to standard output in WAV, AIFF, AIFF-C or raw
format. Currently the -g option is required by the NetBSD
version of cdparanoia. A hypothetical example of how to save
track 2 as a WAV file is as follows:
cdparanoia -g /dev/rcd0d 2 track-02.wav
If you want to grab all files from a CD, cdparanoia's batch mode is useful:
cdparanoia -g /dev/rcd0d -B
For ATAPI or SCSI CD-ROMs the
audio/cdd package can be
used. To extract track 2 with cdd, type:
cdd -t 2 `pwd`
This will put a file called
in the current directory.
For SCSI CD-ROMS the
audio/tosha package can be used.
To extract track 2 with tosha, you should be able to type:
tosha -d-t 2 -o track-02.cda
The data can then be post-processed e.g. by encoding it into MP3 streams (see Section 13.9, “Creating an MP3 (MPEG layer 3) file from an audio CD”) or by writing them to CD-Rs (see Section 13.11, “Using a CD-R writer to create audio CDs”).
The basic steps in creating an MPEG layer 3 (MP3) file from an audio CD (using software from the NetBSD packages collection) are:
Extract (rip) the audio data of the CD as shown in Section 13.8, “Using audio CDs with NetBSD”.
Convert the CD audio format file to WAV format. You only need to perform this job if your ripping program (e.g. tosha, cdd) didn't already do the job for you!
$sox -s -w -c 2 -r 44100 -t cdr track-02.cda track-02.wav
This will convert
in raw CD format to
track-02.wav in WAV format,
using signed 16-bit
words with 2
channels at a sampling
Encode the WAV file into MP3 format.
$bladeenc -128 -QUIT track-02.wav
This will encode
MP3 format, using a bit rate if
for bladeenc describes bit-rates in more detail.
$lame -p -o -v -V 5 -h track-02.wav track-02.mp3
You may wish to use a lower quality, depending on your taste and hardware.
The process of writing a CD consists of two steps: First, a "image" of the data must be generated, which can then be written to CD-R in a second step.
Reading an pre-existing ISO image
dd if=/dev/rcd0a of=filename.iso bs=2k
Alternatively, you can create a new ISO image yourself:
Generating the ISO image
Put all the data you want to put on CD into one directory. Next you need to generate a disk-like ISO image of your data. The image stores the data in the same form as they're later put on CD, using the ISO 9660 format. The basic ISO9660 format only understands 8+3 filenames (max. eight letters for filename, plus three more for an extension). As this is not practical for Unix filenames, a so-called "Rockridge Extension" needs to be employed to get longer filenames. (A different set of such extension exists in the Microsoft world, to get their long filenames right; that's what's known as Joliet filesystem).
The ISO image is created using the mkisofs command,
which is part
Example: if you have your data in /usr/tmp/data, you can generate a ISO image file in /usr/tmp/data.iso with the following command:
$mkisofs -o data.iso -r data Using NETBS000.GZ;1 for data/binary/kernel/netbsd.INSTALL.gz (netbsd.INSTALL_TINY.gz) Using NETBS001.GZ;1 for data/binary/kernel/netbsd.GENERIC.gz (netbsd.GENERIC_TINY.gz) 5.92% done, estimate finish Wed Sep 13 21:28:11 2000 11.83% done, estimate finish Wed Sep 13 21:28:03 2000 17.74% done, estimate finish Wed Sep 13 21:28:00 2000 23.64% done, estimate finish Wed Sep 13 21:28:03 2000 ... 88.64% done, estimate finish Wed Sep 13 21:27:55 2000 94.53% done, estimate finish Wed Sep 13 21:27:55 2000 Total translation table size: 0 Total rockridge attributes bytes: 5395 Total directory bytes: 16384 Path table size(bytes): 110 Max brk space used 153c4 84625 extents written (165 Mb)
Please see the mkisofs(8) man page for other options like noting publisher and preparer. The Bootable CD ROM How-To explains how to generate a bootable CD.
Writing the ISO image to CD-R
When you have the ISO image file,
you just need to write it on a
CD. This is done with the "cdrecord" command from the
Insert a blank CD-R, and off we go:
cdrecord -v dev=/dev/rcd0d data.iso...
After starting the command, 'cdrecord' shows you a lot of information about your drive, the disk and the image you're about to write. It then does a 10 seconds countdown, which is your last chance to stop things - type ^C if you want to abort. If you don't abort, the process will write the whole image to the CD and return with a shell prompt.
Note that cdrecord(8) works on both SCSI and IDE (ATAPI) drives.
Mount the just-written CD and test it as you would do with any "normal" CD, see Section 13.3, “Reading data CDs with NetBSD”.
If you want to make a backup copy of one of your audio CDs, you can do so by extracting ("ripping") the audio tracks from the CD, and then writing them back to a blank CD. Of course this also works fine if you only extract single tracks from various CDs, creating your very own mix CD!
The steps involved are:
If you have converted all your audio CDs to MP3 and now want to make a mixed CD for your (e.g.) your car, you can do so by first converting the .mp3 files back to .wav format, then write them as a normal audio CD.
The steps involved here are:
Create .wav files from your .mp3 files:
mpg123 -w foo.wav foo.mp3
Do this for all of the MP3 files that you want to have on your audio CD. The .wav filenames you use don't matter.
Write the .wav files to CD as described under Section 13.11, “Using a CD-R writer to create audio CDs”.
To copy an audio CD while not introducing any pauses as mandated by the CDDA standard, you can use cdrdao for that:
cdrdao read-cd --device /dev/rcd0d data.toc
cdrdao write --device /dev/rcd1d data.toc
If you have both a CD-R and a CD-ROM drive in your machine, you can copy a data CD with the following command:
cdrecord dev=/dev/rcd1d /dev/rcd0d
Here the CD-ROM (cd0) contains the CD you want to copy, and the CD-R
(cd1) contains the blank disk. Note that this only works with computer
disks that contain some sort of data, it does
not work with
audio CDs! In practice you'll also want to add something like
speed=8" to make things a bit
You can treat a CD-RW drive like a CD-R drive (see Section 13.10, “Using a CD-R writer with data CDs”) in NetBSD, creating images with mkisofs(8) and writing them on a CD-RW medium with cdrecord(8).
If you want to blank a CD-RW, you can do this with cdrecord's
cdrecord dev=/dev/rcd0d blank=fast
There are several other ways to blank the CD-RW,
call cdrecord(8) with
blank=help" for a list. See the cdrecord(8)
man page for more information.
Currently, NetBSD supports DVD media through the ISO 9660
also used for CD-ROMs. The new UDF filesystem also present on DVDs
has been supported since NetBSD 4.0. Information about mounting ISO 9660
and UDF filesystems can be found in the mount_cd9660(8) and
mount_udf(8) manual pages respectively.
DVDs, DivX and many avi files be played with
For some hints on creating DVDs, see this postings about growisofs and this article about recording CDs and DVDs with NetBSD.
To create an ISO image and save the checksum do this:
readcd dev=/dev/cd0d f=/tmp/cd.iso
Here is an alternative using dd(1):
dd if=/dev/cd0d of=/tmp/cd.iso bs=2048
If the CD has errors you can recover the rest with this:
dd if=/dev/cd0d of=/tmp/cd.iso bs=2048 conv=noerror
To create an ISO image from a mounted data CD first, mount the CD disk by:
mount -t cd9660 -r /dev/cd0d /mnt/cdrom
Second, get the image:
mkhybrid -v -l -J -R -o /tmp/my_cd.iso /mnt/cdrom/
You can read the volume data from an unmounted CD with this command:
file -s /dev/cd0d
You can read the volume data from an ISO image with this command:
isoinfo -d -i /tmp/my_cd.iso
You can get the unique disk number from an unmounted CD with this:
You can read the table of contents of an unmounted CD with this command:
cdrecord -v dev=/dev/cd0d -toc |
|This article may be expanded with text translated from the corresponding article in the German Wikipedia. (January 2012)|
||This article needs additional citations for verification. (February 2007)|
||This article needs attention from an expert in Psychology. (October 2009)|
Gestalt psychology or gestaltism (German: Gestalt – "essence or shape of an entity's complete form") is a theory of mind and brain of the Berlin School; the operational principle of gestalt psychology is that the brain is holistic, parallel, and analog, with self-organizing tendencies. The principle maintains that the human eye sees objects in their entirety before perceiving their individual parts, suggesting the whole is greater than the sum of its parts. Gestalt psychology tries to understand the laws of our ability to acquire and maintain stable percepts in a noisy world. Gestalt psychologists stipulate that perception is the product of complex interactions among various stimuli. Contrary to the behaviorist approach to understanding the elements of cognitive processes, gestalt psychologists sought to understand their organization (Carlson and Heth, 2010). The gestalt effect is the form-generating capability of our senses, particularly with respect to the visual recognition of figures and whole forms instead of just a collection of simple lines and curves. In psychology, gestaltism is often opposed to structuralism. The phrase "The whole is greater than the sum of the parts" is often used when explaining gestalt theory, though this is a mistranslation of Kurt Koffka's original phrase, "The whole is other than the sum of the parts". Gestalt theory allows for the breakup of elements from the whole situation into what it really is.
The concept of gestalt was first introduced in contemporary philosophy and psychology by Christian von Ehrenfels (a member of the School of Brentano). The idea of gestalt has its roots in theories by David Hume, Johann Wolfgang von Goethe, Immanuel Kant, David Hartley, and Ernst Mach. Max Wertheimer's unique contribution was to insist that the "gestalt" is perceptually primary, defining the parts it was composed from, rather than being a secondary quality that emerges from those parts, as von Ehrenfels's earlier Gestalt-Qualität had been.
Both von Ehrenfels and Edmund Husserl seem to have been inspired by Mach's work Beiträge zur Analyse der Empfindungen (Contributions to the Analysis of Sensations, 1886), in formulating their very similar concepts of gestalt and figural moment, respectively. On the philosophical foundations of these ideas see Foundations of Gestalt Theory (Smith, ed., 1988).
Early 20th century theorists, such as Kurt Koffka, Max Wertheimer, and Wolfgang Köhler (students of Carl Stumpf) saw objects as perceived within an environment according to all of their elements taken together as a global construct. This 'gestalt' or 'whole form' approach sought to define principles of perception—seemingly innate mental laws that determined the way objects were perceived. It is based on the here and now, and in the way things are seen. Images can be divided into figure or ground. The question is what is perceived at first glance: the figure in front, or the background.
These laws took several forms, such as the grouping of similar, or proximate, objects together, within this global process. Although gestalt has been criticized for being merely descriptive, it has formed the basis of much further research into the perception of patterns and objects (Carlson et al. 2000), and of research into behavior, thinking, problem solving and psychopathology.
Gestalt therapy
The founders of Gestalt therapy, Fritz and Laura Perls, had worked with Kurt Goldstein, a neurologist who had applied principles of Gestalt psychology to the functioning of the organism. Laura Perls had been a Gestalt psychologist before she became a psychoanalyst and before she began developing Gestalt therapy together with Fritz Perls. The extent to which Gestalt psychology influenced Gestalt therapy is disputed, however. In any case it is not identical with Gestalt psychology. On the one hand, Laura Perls preferred not to use the term "Gestalt" to name the emerging new therapy, because she thought that the gestalt psychologists would object to it; on the other hand Fritz and Laura Perls clearly adopted some of Goldstein's work. Thus, though recognizing the historical connection and the influence, most gestalt psychologists emphasize that gestalt therapy is not a form of gestalt psychology.
Theoretical framework and methodology
The investigations developed at the beginning of the 20th century, based on traditional scientific methodology, divided the object of study into a set of elements that could be analyzed separately with the objective of reducing the complexity of this object. Contrary to this methodology, the school of gestalt practiced a series of theoretical and methodological principles that attempted to redefine the approach to psychological research.
The theoretical principles are the following:
- Principle of Totality—The conscious experience must be considered globally (by taking into account all the physical and mental aspects of the individual simultaneously) because the nature of the mind demands that each component be considered as part of a system of dynamic relationships.
- Principle of psychophysical isomorphism – A correlation exists between conscious experience and cerebral activity.
Based on the principles above the following methodological principles are defined:
- Phenomenon experimental analysis—In relation to the Totality Principle any psychological research should take as a starting point phenomena and not be solely focused on sensory qualities.
- Biotic experiment—The school of gestalt established a need to conduct real experiments that sharply contrasted with and opposed classic laboratory experiments. This signified experimenting in natural situations, developed in real conditions, in which it would be possible to reproduce, with higher fidelity, what would be habitual for a subject.
Support from cybernetics and neurology
In the 1940s and 1950s, laboratory research in neurology and what became known as cybernetics on the mechanism of frogs' eyes indicate that perception of 'gestalts' (in particular gestalts in motion) is perhaps more primitive and fundamental than 'seeing' as such:
- A frog hunts on land by vision... He has no fovea, or region of greatest acuity in vision, upon which he must center a part of the image... The frog does not seem to see or, at any rate, is not concerned with the detail of stationary parts of the world around him. He will starve to death surrounded by food if it is not moving. His choice of food is determined only by size and movement. He will leap to capture any object the size of an insect or worm, providing it moves like one. He can be fooled easily not only by a piece of dangled meat but by any moving small object... He does remember a moving thing provided it stays within his field of vision and he is not distracted. Cyberneticist Valentin Turchin points out that the gestalts observed in what we usually imagine are 'still images' are exactly the kind of 'moving objects' that make the frog's retina respond:
- The lowest-level concepts related to visual perception for a human being probably differ little from the concepts of a frog. In any case, the structure of the retina in mammals and in human beings is the same as in amphibians. The phenomenon of distortion of perception of an image stabilized on the retina gives some idea of the concepts of the subsequent levels of the hierarchy. This is a very interesting phenomenon. When a person looks at an immobile object, "fixes" it with his eyes, the eyeballs do not remain absolutely immobile; they make small involuntary movements. As a result the image of the object on the retina is constantly in motion, slowly drifting and jumping back to the point of maximum sensitivity. The image "marks time" in the vicinity of this point.
Emergence is the process of complex pattern formation from simpler rules. It is demonstrated by the perception of the dog picture, which depicts a Dalmatian dog sniffing the ground in the shade of overhanging trees. The dog is not recognized by first identifying its parts (feet, ears, nose, tail, etc.), and then inferring the dog from those component parts. Instead, the dog is perceived as a whole, all at once. However, this is a description of what occurs in vision and not an explanation. Gestalt theory does not explain how the percept of a dog emerges.
Reification is the constructive or generative aspect of perception, by which the experienced percept contains more explicit spatial information than the sensory stimulus on which it is based.
For instance, a triangle is perceived in picture A, though no triangle is there. In pictures B and D the eye recognizes disparate shapes as "belonging" to a single shape, in C a complete three-dimensional shape is seen, where in actuality no such thing is drawn.
Reification can be explained by progress in the study of illusory contours, which are treated by the visual system as "real" contours.
Multistability (or multistable perception) is the tendency of ambiguous perceptual experiences to pop back and forth unstably between two or more alternative interpretations. This is seen for example in the Necker cube, and in Rubin's Figure/Vase illusion shown here. Other examples include the Three-legged blivet and artist M. C. Escher's artwork and the appearance of flashing marquee lights moving first one direction and then suddenly the other. Again, gestalt does not explain how images appear multistable, only that they do.
Invariance is the property of perception whereby simple geometrical objects are recognized independent of rotation, translation, and scale; as well as several other variations such as elastic deformations, different lighting, and different component features. For example, the objects in A in the figure are all immediately recognized as the same basic shape, which are immediately distinguishable from the forms in B. They are even recognized despite perspective and elastic deformations as in C, and when depicted using different graphic elements as in D. Computational theories of vision, such as those by David Marr, have had more success in explaining how objects are classified.
Emergence, reification, multistability, and invariance are not necessarily separable modules to model individually, but they could be different aspects of a single unified dynamic mechanism.
The fundamental principle of gestalt perception is the law of prägnanz (in the German language, pithiness), which says that we tend to order our experience in a manner that is regular, orderly, symmetric, and simple. Gestalt psychologists attempt to discover refinements of the law of prägnanz, and this involves writing down laws that, hypothetically, allow us to predict the interpretation of sensation, what are often called "gestalt laws". These include:
Gestalt laws of grouping
A major aspect of Gestalt psychology is that it implies that the mind understands external stimuli as whole rather than the sum of their parts. The wholes are structured and organized using grouping laws. The various laws are called laws or principles, depending on the paper where they appear—but for simplicity sake, this article uses the term laws. These laws deal with the sensory modality vision however there are analogous laws for other sensory modalities including auditory, tactile, gustatory and olfactory (Bregman – GP). The visual Gestalt principles of grouping were introduced in Wertheimer (1923). Through the 1930s and '40s Wertheimer, Kohler and Koffka formulated many of the laws of grouping through the study of visual perception.
Law of Proximity—The law of proximity states that when an individual perceives an assortment of objects they perceive objects that are close to each other as forming a group. For example, in the figure that illustrates the Law of proximity, there are 72 circles, but we perceive the collection of circles in groups. Specifically, we perceive there is a group of 36 circles on the left side of the image, and three groups of 12 circles on the right side of the image. This law is often used in advertising logos to emphasize which aspects of events are associated.
Law of Similarity—The law of similarity states that elements within an assortment of objects are perceptually grouped together if they are similar to each other. This similarity can occur in the form of shape, colour, shading or other qualities. For example, the figure illustrating the law of similarity portrays 36 circles all equal distance apart from one another forming a square. In this depiction, 18 of the circles are shaded dark and 18 of the circles are shaded light. We perceive the dark circles as grouped together, and the light circles as grouped together forming six horizontal lines within the square of circles. This perception of lines is due to the law of similarity.
Law of Closure—The law of closure states that individuals perceive objects such as shapes, letters, pictures, etc., as being whole when they are not complete. Specifically, when parts of a whole picture are missing, our perception fills in the visual gap. Research shows that the reason the mind completes a regular figure that is not perceived through sensation is to increase the regularity of surrounding stimuli. For example, the figure that depicts the law of closure portrays what we perceive as a circle on the left side of the image and a rectangle on the right side of the image. However, gaps are missing from the shapes. If the law of closure did not exist, the image would depict an assortment of different lines with different lengths, rotations, and curvatures—but with the law of closure, we perceptually combine the lines into whole shapes.
Law of Symmetry—The law of symmetry states that the mind perceives objects as being symmetrical and forming around a center point. It is perceptually pleasing to divide objects into an even number of symmetrical parts. Therefore, when two symmetrical elements are unconnected the mind perceptually connects them to form a coherent shape. Similarities between symmetrical objects increase the likelihood that objects are grouped to form a combined symmetrical object. For example, the figure depicting the law of symmetry shows a configuration of square and curled brackets. When the image is perceived, we tend to observe three pairs of symmetrical brackets rather than six individual brackets.
Law of Common Fate—The law of common fate states that objects are perceived as lines that move along the smoothest path. Experiments using the visual sensory modality found that movement of elements of an object produce paths that individuals perceive that the objects are on. We perceive elements of objects to have trends of motion, which indicate the path that the object is on. The law of continuity implies the grouping together of objects that have the same trend of motion and are therefore on the same path. For example, if there are an array of dots and half the dots are moving upward while the other half are moving downward, we would perceive the upward moving dots and the downward moving dots as two distinct units.
Law of Continuity—The law of continuity states that elements of objects tend to be grouped together, and therefore integrated into perceptual wholes if they are aligned within an object. In cases where there is an intersection between objects, individuals tend to perceive the two objects as two single uninterrupted entities. Stimuli remain distinct even with overlap. We are less likely to group elements with sharp abrupt directional changes as being one object.
Law of Good Gestalt—The law of good gestalt explains that elements of objects tend to be perceptually grouped together if they form a pattern that is regular, simple, and orderly. This law implies that as individuals perceive the world, they eliminate complexity and unfamiliarity so they can observe a reality in its most simplistic form. Eliminating extraneous stimuli helps the mind create meaning. This meaning created by perception implies a global regularity, which is often mentally prioritized over spatial relations. The law of good gestalt focuses on the idea of conciseness, which is what all of gestalt theory is based on. This law has also been called the law of Prägnanz. Prägnanz is a German word that directly translates to mean "pithiness" and implies the ideas of salience, conciseness and orderliness.
Law of Past Experience—The law of past experience implies that under some circumstances visual stimuli are categorized according to past experience. If two objects tend to be observed within close proximity, or small temporal intervals, the objects are more likely to be perceived together. For example, the English language contains 26 letters that are grouped to form words using a set of rules. If an individual reads an English word they have never seen, they use the law of past experience to interpret the letters "L" and "I" as two letters beside each other, rather than using the law of closure to combine the letters and interpret the object as an uppercase U.
The gestalt laws of grouping have recently been subjected to modern methods of scientific evaluation by examining the visual cortex using cortical algorithms. Current Gestalt psychologists have described their findings, which showed correlations between physical visual representations of objects and self-report perception as the laws of seeing.
Gestalt views in psychology
Gestalt psychologists find it is important to think of problems as a whole. Max Wertheimer considered thinking to happen in two ways: productive and reproductive.
Productive thinking is solving a problem with insight.
This is a quick insightful unplanned response to situations and environmental interaction.
Reproductive thinking is solving a problem with previous experiences and what is already known. (1945/1959).
This is a very common thinking. For example, when a person is given several segments of information, he/she deliberately examines the relationships among its parts, analyzes their purpose, concept, and totality, he/she reaches the "aha!" moment, using what is already known. Understanding in this case happens intentionally by reproductive thinking.
Another gestalt psychologist, Perkins, believes insight deals with three processes:
- Unconscious leap in thinking.
- The increased amount of speed in mental processing.
- The amount of short-circuiting that occurs in normal reasoning.
Views going against the gestalt psychology are:
- Nothing-special view
- Neo-gestalt view
- The Three-Process View
Gestalt psychology should not be confused with the gestalt therapy of Fritz Perls, which is only peripherally linked to gestalt psychology. A strictly gestalt psychology-based therapeutic method is Gestalt Theoretical Psychotherapy, developed by the German gestalt psychologist and psychotherapist Hans-Jürgen Walter.
Fuzzy-trace theory
Fuzzy-trace theory, a dual process model of memory and reasoning, was also derived from Gestalt Psychology. Fuzzy-trace theory posits that we encode information into two separate traces: verbatim and gist. Information stored in verbatim is exact memory for detail (the individual parts of a pattern, for example) while information stored in gist is semantic and conceptual (what we perceive the pattern to be). The effects seen in Gestalt psychology can be attributed to the way we encode information as gist.
Uses in human–computer interaction
The gestalt laws are used in user interface design. The laws of similarity and proximity can, for example, be used as guides for placing radio buttons. They may also be used in designing computers and software for more intuitive human use. Examples include the design and layout of a desktop's shortcuts in rows and columns. Gestalt psychology also has applications in computer vision for trying to make computers "see" the same things as humans do.
Quantum cognition modeling
Similarities between Gestalt phenomena and quantum mechanics have been pointed out by, among others, chemist Anton Amann, who commented that "similarities between Gestalt perception and quantum mechanics are on a level of a parable" yet may give useful insight nonetheless. Physicist Elio Conte and co-workers have proposed abstract, mathematical models to describe the time dynamics of cognitive associations with mathematical tools borrowed from quantum mechanics and has discussed psychology experiments in this context. A similar approach has been suggested by physicists David Bohm, Basil Hiley and philosopher Paavo Pylkkänen with the notion that mind and matter both emerge from an "implicate order". The models involve non-commutative mathematics; such models account for situations in which the outcome of two measurements performed one after the other can depend on the order in which they are performed—a pertinent feature for psychological processes, as it is obvious that an experiment performed on a conscious person may influencing the outcome of a subsequent experiment by changing the state of mind of that person.
In some scholarly communities, such as cognitive psychology and computational neuroscience, gestalt theories of perception are criticized for being descriptive rather than explanatory in nature. For this reason, they are viewed by some as redundant or uninformative. For example, Bruce, Green & Georgeson conclude the following regarding gestalt theory's influence on the study of visual perception:
- The physiological theory of the gestaltists has fallen by the wayside, leaving us with a set of descriptive principles, but without a model of perceptual processing. Indeed, some of their "laws" of perceptual organisation today sound vague and inadequate. What is meant by a "good" or "simple" shape, for example?
See also
- Gestalt therapy—often mistaken for gestalt psychology
- Structural information theory
- Rudolf Arnheim
- Wolfgang Metzger
- Kurt Goldstein
- Pál Schiller Harkai
- Solomon Asch
- Hermann Friedmann
- James J. Gibson
- James Tenney
- Graz School
- Important publications in gestalt psychology
- Optical illusion
- Pattern recognition (psychology)
- Pattern recognition (machine learning)
- Amodal perception
- Topological data analysis
- Fuzzy-trace theory
- Laws of Association
- David Hothersall: History of Psychology, chapter seven,(2004)
- Tuck, Michael. "Gestalt Principles Applied in Design (Aug 17 2010)". Retrieved 11/12/11.
- Humphrey, G. (1924). The psychology of the gestalt. Journal of Educational Psychology, 15(7), 401–412. doi:10.1037/h0070207
- Bernd Bocian:Fritz Perls in Berlin 1893–1933. Expressionism – Psychonalysis – Judaism, 2010, p. 190, EHP Verlag Andreas Kohlhage, Bergisch Gladbach.
- Joe Wysong/Edward Rosenfeld (eds): An Oral History of Gestalt Therapy, Highland, New York 1982, The Gestalt Journal Press, p. 12.
- Allen R. Barlow, "Gestalt-Antecedent Influence or Historical Accident", The Gestalt Journal, Volume IV, Number 2, (Fall, 1981)
- Mary Henle noted in her presidential address to Division 24 at the meeting of the American Psychological Association (1975): "What Perls has done has been to take a few terms from Gestalt psychology, stretch their meaning beyond recognition, mix them with notions—often unclear and often incompatible —– from the depth psychologies, existentialism, and common sense, and he has called the whole mixture gestalt therapy. His work has no substantive relation to scientific Gestalt psychology. To use his own language, Fritz Perls has done 'his thing'; whatever it is, it is not Gestalt psychology". Gestalt theory. However she restricts herself explicitly to only three of Perls' books from 1969 and 1972, leaving out Perls' earlier work, and Gestalt therapy in general. See Barlow criticizing Henle: Allen R. Barlow: Gestalt Therapy and Gestalt Psychology. Gestalt – Antecedent Influence or Historical Accident, in: The Gestalt Journal, Volume IV, Number 2, Fall, 1981.
- William Ray Woodward, Robert Sonné Cohen – World views and scientific discipline formation: science studies in the German Democratic Republic : papers from a German-American summer institute, 1988
- Lettvin, J.Y., Maturana, H.R., Pitts, W.H., and McCulloch, W.S. (1961). Two Remarks on the Visual System of the Frog. In Sensory Communication edited by Walter Rosenblith, MIT Press and John Wiley and Sons: New York
- Valentin Fedorovich Turchin – The phenomenon of science – a cybernetic approach to human evolution – Columbia University Press, 1977
- "Gestalt Isomorphism". Sharp.bu.edu. Retrieved 2012-04-06.
- Sternberg, Robert, Cognitive Psychology Third Edition, Thomson Wadsworth© 2003.
- Stevenson, Herb. "Emergence: The Gestalt Approach to Change". Unleashing Executive and Orzanizational Potential. Retrieved 7 April 2012. Text "noedit" ignored (help)
- Soegaard, Mads. "Gestalt Principles of form Perception". Interaction Design. Retrieved 8 April 2012. Text "noedit" ignored (help)
- Todorovic, Dejan. "Gestalt Principles". scholarpedia. Retrieved 5 April 2012. Text "noedit" ignored (help)
- Langley& associates, 1987; Perkins, 1981; Weisberg, 1986,1995">
- Reyna, Valerie (2012). "A new institutionism: Meaning, memory, and development in Fuzzy-Trace Theory". Judgment and Decision Making 7 (3): 332–359.
- Soegaard, Mads. "Gestalt principles of form perception". Interaction-design.org. Retrieved 2012-04-06.
- Elio Conte, Orlando Todarello, Antonio Federici, Francesco Vitiello, Michele Lopane, Andrei Khrennikov, Joseph P. Zbilut: Some remarks on an experiment suggesting quantum-like behavior of cognitive entities and formulation of an abstract quantum mechanical formalism to describe cognitive entity and its dynamics, Chaos, Solitons & Fractals, vol. 31, no. 5, March 2007, pp. 1076–1088 doi:10.1016/j.chaos.2005.09.061, arXiv:0710.5092 (submitted 26 October 2007)
- Elio Conte, Orlando Todarello, Antonio Federici, Francesco Vitiello, Michele Lopane, Andrei Khrennikov: A Preliminar Evidence of Quantum Like Behavior in Measurements of Mental States, arXiv:quant-ph/0307201 (submitted 28 July 2003)
- B.J. Hiley: Particles, fields, and observers, Volume I The Origins of Life, Part 1 Origin and Evolution of Life, Section II The Physical and Chemical Basis of Life, pp. 87–106 (PDF)
- Basil J. Hiley, Paavo Pylkkänen: Naturalizing the mind in a quantum framework. In Paavo Pylkkänen and Tere Vadén (eds.): Dimensions of conscious experience, Advances in Consciousness Research, Volume 37, John Benjamins B.V., 2001, ISBN 90-272-5157, pages 119-144
- Bruce, V., Green, P. & Georgeson, M. (1996). Visual perception: Physiology, psychology and ecology (3rd ed.). LEA. p. 110.
- Carlson, Neil R. and Heth, C. Donald (2010) Psychology the Science of Behaviour Ontario, CA: Pearson Education Canada. pp 20–22.
- Smith, Barry (ed.) (1988) Foundations of Gestalt Theory, Munich and Vienna: Philosophia Verlag, 1988.
- Gestalt psychology on Encyclopædia Britannica
- Gestalt Society of Croatia
- International Society for Gestalt Theory and its Applications – GTA
- Embedded Figures in Art, Architecture and Design
- On Max Wertheimer and Pablo Picasso
- On Esthetics and Gestalt Theory
- The World In Your Head – by Steven Lehar
- Gestalt Isomorphism and the Primacy of Subjective Conscious Experience – by Steven Lehar
- The new gestalt psychology of the 21st century
- The Pennsylvania Gestalt Center
- Gestalt Theory
- Ecological Approach to Visual Perception
- James J. Gibson in brief |
ixated as the Victorians were on death, it is not at all surprising that a number of talented British poets in the nineteenth century would have explored mourning. From the elaborate, regal funerals of British royalty down to the simpler funerals of members of the working class — paid for by dutifully purchased death insurance — rituals of death formed an integral part of British life (Jones, 199). Yet the poets explored here do not write about funerals. Nor do they follow the traditional patterns of elegiac poetry, for example, used in Algernon Charles Swinburne's elegy to Baudelaire, "Ave Atque Vale." Rather, they thoroughly explore various aspects of the emotional experience of mourning. Some poets, notably Elizabeth Barrett Browning and Alfred Lord Tennyson, delve into the ways that people express grief. Tennyson also struggles with the role of one's conscience in the process of mourning and with ways to reconcile a personal need for comfort and consolation with a strong desire to honor and preserve the memory of the deceased. Christina Rossetti examines how gender and sexual interest influence the way people mourn; in particular, she considers the value of mourning in relationships in which the fleeting nature of love is acknowledged. Both she and Tennyson deal with the dead woman as an aesthetic object, and the belated mourning of a previously inattentive male figure who feasts his eyes upon the dead female, bringing to light another aspect of the way gender relations affect mourning. Finally, Rossetti presents contrasting images of the afterlife that variously affect the practice of mourning, depending upon the consciousness of the deceased.
A common motif permeating poetry that deals with mourning is sound, be it in the form of tears of mourning, a missed language, a song of mourning or a noted silence. A large number of poems that consider death and mourning utilize images related to sound. The role of sound in the various elements of mourning these poets consider is crucial: sound is a critical link between two people; when one person dies, this connection is apparently broken, and the mourner is left in the unhappy position of vocally trying to maintain contact with the deceased. Thus, the numerous sound illusions in elegiac poetry are no more surprising than the general Victorian focus on death: people have mourned by means of words and songs for a long time. The emotionally charged and experimental ways of examining this method of mourning set apart the work of these poets.
Seeking a means of expressing grief
The loss of her parents is a shadow that hangs over Aurora Leigh in her struggle to mold her identity, in Elizabeth Barrett Browning's Aurora Leigh. Aurora tells of the loss of her mother early in the poem, describing the pain of a young child hardly old enough to understand what has happened but intensely aware that something is amiss. Aurora feels that she will always be searching for a mother because she was so young when her mother died and hence their time together was extremely brief.
I felt a mother-want about the world,
And still went seeking, like a bleating lamb
Left out at night, in shutting up the fold, —
As restless as a nest-deserted bird
Grown chill through something being away, though what
It knows not. [First book, lines 40-45, p. 2]
In Aurora's pain, she likens herself to "a bleating lamb," interestingly connecting her own mourning with the repeated cry of a young animal. The other animal analogy she makes is to a "nest-deserted bird," and since birds are associated with songs, this also seems to have aural connotations. Aurora is like a crying animal, calling out for attention to a loved one who will not hear her. Thus, this passage subtly develops the idea of mourning as a plea to the one who has died somehow to respond and return. Aurora feels that her mother, by dying, has shut her out of her love and deserted her.
The animals symbolically carry out Aurora's crying; she is too young to mourn by means of words, tears or songs, as her father writes in the verses to the memory of her mother, "Weep for an infant too young to weep much / When death removed this mother" (first book, lines 103-4, p. 4). Her father's entreaty is primarily calling for sympathy for the child who is too young to have known her mother and therefore too young to mourn her. However, he is literally calling for the act of mourning by means of tears to make up for the child's inability to cry about death at her age. Thus, there is a positive association with mourning in Aurora Leigh: there is a perceived need for Aurora to cry as part of the process of mourning her mother; the sadness of this circumstance is that the child is too young to mourn and hence show reverence and the pain of loss. She is also too young to experience the cathartic effects of mourning.
When she is still quite yo ung, Aurora's father dies as well. Aurora is then taken from her nurse and from Italy, her homeland, and sent to live with her aunt. The moment of parting from her nurse is wrought with agony and sounds of loss. Aurora cannot yet express her sadness, but at this point she is acutely aware of others' sounds, particularly the sounds emitted by her nurse in her unhappiness at being parted from Aurora.
[W]ith a shriek,
She let me go, — while I, with ears too full
Of my father's silence, to shriek back a word,
In all a child's astonishment at grief
Stared at the wharfage where she stood and moaned!
My poor Assunta, where she stood and moaned! [First Book, lines 226-31, p. 7]
Even as Aurora cannot speak in her youthful misery, she notices that her father's voice has been stilled and that her nurse moans as she anticipates the loss of Aurora. Thus, vocal mourning is associated not only with the adult way of grieving a death but also with the separation of people. Sound creates a connection between people, and death and parting break that link.
It is not until Aurora reaches England that she is finally able to cry over the death of her father. She hears Britons speaking English, a language which she had only heard her father speak, and she feels the pain of unfulfilled expectations: the love she had come to associate with that language has passed with the death of her father, and that same language is now spoken by strangers.
And when I heard my father's language first
From alien lips which had no kiss for mine,
I wept aloud, then laughed, then wept, then wept, —
And some one near me said the child was mad
Through much sea-sickness. [First Book, lines 254-58, p. 8]
Hearing the language and focusing on her loss, Aurora is finally able to grieve aloud, and she cries and laughs with sadness and relief at her ability to mourn. Her expression of grief is perceived by those around her as delirium or a fit of insanity. Although this is a somewhat negative response to her mourning, those around her do not know the cause of her tears, and so this passage does not seem to contradict the general sentiment in Aurora Leigh with regards to vocal mourning: it is a natural expression of pain that adults manifest, and that children, as they mature, will also learn to convey. Furthermore, the loss of a person creates a void: his voice is silenced. Hence, the vocal act of calling a person back fills that sound void, and is therefore comforting to the mourner.
Tennyson also considers the appropriate way to mourn and the role of sound in this process — particularly in the form of mourning songs — in In Memoriam, his experimental elegy for his beloved friend, Arthur Henry Hallam. He questions the morality of such methods of mourning with respect to the memory of the deceased; yet it seems that in addition to his concern for the deceased, the narrator is also personally dissatisfied with the effects of mourning songs.
Peace; come away: the song of woe
Is after all earthly song.
Peace; come away: we do him wrong
To sing so wildly: let us go. [57, lines 1-4, p. 239]
These lines suggest that the songs a mourner sings and the sounds of weeping together express a need to cling to the deceased and a refusal to relinquish him. The speaker and the other mourners dwell on the dead one, carrying on wildly rather than letting him pass peacefully and calmly. The narrator's call to those singing is to leave the dead behind and go on. He literally says that it does the deceased wrong when those alive sing, but he seems to be concerned with those who survive as well. Although this singing about loss is an earthly thing to do, and it is not exactly viewed by the speaker as an inappropriate way for those on earth to channel their grief, it is nonetheless not satisfying for exactly the reason it is natural: it is only earthly, and it provides no divine satisfaction to the mourner. The narrator is seeking that sort of spiritual experience and confirmation that Hallam died for a divine reason. The speaker feels this divine connection to Hallam in poem 95 of In Memoriam, and it is by means of touch not sound that the spiritual experience occurs, and that the narrator finds at least some degree of satisfaction.
The role of conscience
Throughout In Memoriam, Tennyson, as an evolving narrator, struggles in his efforts to cope with Hallam's death. He considers thoroughly how he will be affected by different modes of mourning, and he worries deeply about the effects these will have on him. He is acutely concerned that in mourning Hallam, he will somehow ease his own suffering; although he is urgently searching for divine meaning behind Hallam's death, his conscience continuously discourages him from mourning in any way that consoles him, and this makes his mourning process even more difficult and painful.
He contemplates the morality of writing about his friend's death, fearing that the action of writing, which may mitigate his pain — "Like dull narcotics, numbing pain" (5, line 8, p. 208), is unethical. The act of forming words out of his pain, he fears, will misrepresent his feelings because it is impossible to perfectly capture the state of his soul (5, lines 1-4, p. 208). Similarly, he worries that beyond simply soothing his pain, he is actually using the death of his friend as inspiration for his own artistic creation, just as the yew tree receives its nourishment from the dead bodies buried in graveyards (2, lines 1-4, p. 207); this is a highly disturbing thought to a character so plagued by his conscience. He also worries that the elegy he is writing for Hallam will not provide an enduring tribute to his friend.
Come; let us go: your cheeks are pale;
But half my life I leave behind.
Methinks my friend is richly shrined;
But I shall pass, my work will fail. [57, lines 5-8, p. 239]
Thus, in his mourning, the narrator is aiming not only to grapple with Hallam's death, but to create a worthy memorial for him.
Although the narrator calls for an end to mournful singing, he expresses the belief that he will always be haunted by sounds reminding him of the death of Hallam. He notes that until his hearing fails him or until he himself dies, he shall hear a slow, constant bell announcing the death of his friend repeatedly in his own ears. He also describes hearing the repeated farewells said to those who are dead.
Yet in these ears, till hearing dies,
One set slow bell will seem to toll
The passing of the sweetest soul
That ever look'd with human eyes.
I hear it now, and o'er and o'er,
Eternal greetings to the dead;
And "Ave, Ave, Ave," said,
"Adieu, adieu," for evermore. [57, lines 9-16, p. 239]
Thus, even when the mourner stops his vocal mourning, he nonetheless remains doomed to hear reminders of the death. It is as though the connection between the two men is so strong that not even death can break it, and sounds still draws them to one another.
At this point in the poem, the narrator seems to believe that it is the responsibility of the mourner forever to hear morbid, painful sounds. Interestingly, if memories of the dead can aurally infiltrate the being of the mourner but the mourner must remain silent and experience pain fully, the natural order of the world becomes inverted: rather than the dead being silenced and those who remain alive retaining the power to speak and sing, the reverse essentially is condoned. That is, the narrator may continue to speak and sing, orally expressing his grief, but such vocal manifestations of sadness are not encouraged. On the contrary, the bells prompting him to grieve, which the narrator hears in his mind, are considered righteous and appropriate.
Whereas some poems denigrate the act of mourning aloud, suggesting that it is inappropriate or inadequate, others discourage such forms of mourning in the particular context of romantic relationships and gender relations. For example, in Christina Rossetti's "Song," the narrator's request that the person to whom the poem is addressed not mourn is more a reflection of the dynamic between the characters and an indicator of gender relations than a suggestion of a general dissatisfaction of the narrator with vocal mourning.
When I am dead, my dearest,
Sing no sad songs for me;
Plant thou no roses at my head,
Nor shady cypress tree:
Be the green grass above me
With showers and dewdrops wet;
And if thou wilt remember,
And if thou wilt forget. [lines 1-8, p. 198]
As George Landow has shown in "The Dead Woman Talks Back: Christina Rossetti's Ironic Intonation of the Dead Fair Maiden," Rossetti is toying with and upsetting common ideas of Victorian femininity. The last two lines of the quotation above seem to conform to standard stereotypes of the self-sacrificing female: she does not want to inconvenience her lover or cause him distress by demanding that he mourn her and always remember her. Yet the final lines of the poem, "Haply I may remember, / And haply may forget" (lines 15-16) reverse the previous vision of self-sacrifice and eternal love on the part of the female speaker, suggesting instead that she is content to or at least quite capable of forgetting her lover and that he should do the same.
This reading allows for further consideration of the role of the specific acts of mourning that are mentioned in the poem. As the poem continues, the speaker describes a senseless afterlife in which she will not see, feel or hear anything, including the nightingale's song. Thus, her death will distinctly separate her from the world, and not even the mourning song of the nightingale will reach her. The removal of the speaker from the living world will render her incapable of hearing songs of mourning, and therefore such mourning rituals become ridiculous. In the course of the poem, the narrator gradually challenges the notion that there is eternal love between herself and the speaker. Thus, death can be seen as the appropriate moment at which to accept the rupture of their relationship, since it would not continue forever anyway; there is no need, then, to prolong the connection by means of songs and other customs of mourning. Rossetti boldly states what many people dare not admit: that the person who has died will not benefit from the mourning practices of those who survive him and hence, such customs only serve those who live. In a situation in which the lover is not likely to go on loving the woman who has died, Rossetti's narrator urges him not to mourn in false, showy ways: both of them will soon forget each other anyway.
Beyond projecting herself into the grave, the speaker in Rossetti's "After Death" is already dead and narrates the poem from the grave. She is thus a standard aesthetic object, deemed a beautiful and rich source for poetry (Landow). In this poem, the lamenting words of the dead woman's beloved reach her ears, even though he does not know it.
He leaned above me, thinking that I slept
And could not hear him; but I heard him say:
"Poor Child, poor child:" and as he turned away
Came a deep silence, and I knew he wept. [lines 5-8, p. 200]
In this case, the man's grieving words do provide comfort to the narrator, whom he did not love in her life and whom he can only pity in death. In the silence, she indicates her certainty that he is crying for her. Thus, words and tears of mourning have practical significance in this poem because they can be heard and felt by the deceased. Here the acts of mourning do not fall on deaf ears, and they are conveyed to the reader by the character with the most personal investment in them: the dead woman for whom they were intended.
The dead woman as an aesthetic object also appears at the end of Tennyson's "The Lady of Shalott" when the knight Lancelot notices the beauty of the dead Lady of Shalott. From the moment that she places herself in the boat and sets out down the river, sound plays a key role in setting the scene of her death. As she floats, "the noises of the night"(line 139) are mentioned, but it is the music that she herself creates that is the more powerful sound.
They heard her singing her last song,
The Lady of Shalott.
Heard a carol, mournful, holy,
Chanted loudly, chanted lowly,
Till her blood was frozen slowly,
And her eyes were darkened wholly,
Turned to tower'd Camelot.
For ere she reach'd upon the tide
The first house by the water-side,
Singing in her song she died,
The Lady of Shalott. [lines 143-53, p. 44]
It is as though the Lady of Shalott sounds her own death knell with her song. Her music is her final and only direct interaction with the world, distant as it is, before her death. She has emerged from her bower and sings a song that reaches the ears of those her boat passes on the journey to Camelot. Then, after death, she physically confronts the outside world. Previously, she saw shadows and reflections of the world outside and made no art that the outside world could appreciate. The sound she emits before her death, however, establishes a brief and melancholic connection between her and the world.
Yet when she reaches Camelot, she is "silent" (line 158, p. 44), and her deathly presence causes the merriment in the palace to cease, killing the sounds within (lines 164-5, p. 45). The silence of those at the palace stresses the absence of vocal mourning, emphasizing the distance between her and them, and their inability to feel deep compassion for her. The connection she made with them with her song is largely one-sided: most of them do not respond. Lancelot, however, does respond, but his comments about the beauty of her face and his request for God's blessing on her ring of transience: in death she has caught his fancy for a moment, but he will soon move on and forget her. In the same way that Rossetti argues against the immortality of love in "Song," Tennyson seems to suggest the inherently ephemeral nature of vocal mourning at the end of "The Lady of Shalott." The use of the word "mournful" anticipates her death, and perhaps compensates for the fact that she will not be mourned by anyone who knew her, except in her own self-pitying song.
The post-mortem mourning of the men who did not love or notice the women in life has a dramatic and ironic, yet somewhat shallow effect. In both "After Death" and "The Lady of Shalott," the reader is left with a slight feeling of satisfaction: at least the man has finally noticed the woman, if only for a moment and after her death. An aura of revenge, if only mild revenge, is thereby achieved: the punishment for having ignored the women in life is that the men are doomed to see beauty in the women only when they are already out of reach.
Images of the Afterlife
Rossetti's imaginings of life after death range from the complete alertness seen in "After Death" to the partially aware state of the woman in "Dream Land" to the completely senseless states of the women in "Song" and "Rest." The woman discussed in "Rest" is completely shielded from the sounds of earthly life. However, the sounds discussed in the poem are not sounds of mourning, but rather the sounds that make up daily life. The omniscient narrator calls for the earth to physically isolate and protect the buried woman from "mirth/ With its harsh laughter" (lines 3-4, p. 200) and the "sound of sighs" (line 4, p. 200). The narrator, who several times alludes to sound, claims that the woman will be spared all noise in her rest, and that the silence of death will be "more musical than any song" (line 10, p. 200). The afterlife in this poem resembles that described in "Song" in that the woman escapes to a senseless, timeless place of peace. Although "Rest" does not deal directly with mourning, it reinforces the idea that sound establishes a connection between people; in the woman's desire to escape the world in death, she is entirely isolated, both from people and sound.
Interestingly, in Rossetti's "Dream Land," she creates a post-mortem state somewhere between the full awareness of the narrator in "After Death" and the senseless afterlife described in "Song" and "Rest." The third person narrator in "Dream Land" never states that the woman in the poem is dead, saying rather that she is asleep, which is also implied by the title. However a number of illusions in the poem suggest that her sleep is actually death, or at least strongly resembles such a state: "a perfect rest" (line 17), "rest, for evermore" (line 25), and
Till time shall cease:
Sleep that no pain shall wake;
Night that no morn shall break
Till joy shall overtake
Her perfect peace. [lines 28-32, p. 199].
The final two lines in particular suggest the prospect of heaven after a deathly respite. Although Rossetti describes how the speaker cannot see grain or feel rain, she nonetheless is able to see "the sky look pale" (line 14) and hear "the nightingale / That sadly sings" (lines 15-6). Her awareness of these sights and sounds suggests an increased awareness as opposed to the senseless afterlife anticipated by the narrator of "Song." In "Dream Land" there is no indictment of mourning. Rather, there is a slight indication that sounds can pass between the worlds of the living and the dead or at least deeply sleeping; therein may lie a justification for the nightingale's song or a lover's cries of mourning, despite the indications in "Song" that such methods of mourning are futile.
As a contrast to Rossetti's quieter imaginings of the afterlife, and a parallel to "After Death," it seems reasonable to mention the afterlife envisioned by Robert Browning in his dramatic monologue, "The Bishop Orders his Tomb at Saint Praxed's Church." In "A Strangely Literal Afterlife," I discuss the lack of faith suggested by the fully aware, materialistic, and grudge-holding afterlife the Bishop envisions for himself. Yet despite Rossetti's devout religious tendencies, her depiction of the afterlife, in the poems discussed above, seems less religiously than emotionally focused: she is concerned with the women's connections or lack thereof to the world and people around them, and how death and awareness of sound affect these connections.
Mourning and the connections people hope to make by emitting sounds are relatively specific themes, yet they illuminate a distinct perspective on the broader topics of personal expression, conscience, gender and the afterlife, all subjects explored by the Victorian poets here considered. They created characters who aspire to connect with other people, sometimes those who have died and are therefore out of reach. These characters struggle to adequately express their emotions. In situations of loss, this need to express oneself is heightened as a result of suffering. As they grieve and search for meaning or solace in the wake of a death, some characters feel the strains of a strong conscience, causing them to question which forms of mourning and self-expression are moral and which are not. Certain poets considered the way that men grieve for women, sometimes viewing the process from unusual angles. Gender relations played an ever-changing role in Victorian life, so of course their appearance in the elegiac poetry of the nineteenth century is understandable. Poetry that deals with love often also considers loss, and this loss allows for creative conceptions of the afterlife. All of these subjects intimately relate to human relationships. By addressing these issues, the poems discussed show that sounds both draw people together and ultimately force them to accept that their connections have been severed — and that all that remains is silence.
Bolton, John, Robert Glorney and Julia Bolton Holloway, eds. Elizabeth Barrett Browning: Aurora Leigh and other Poems. Penguin Putnam Inc., New York, 1995.
Buckley, Jerome H. ed. The Pre-Raphaelites. Academy Chicago Publishers, Chicago: 1968.
Hill, Robert W. Jr., ed. Tennyson's Poetry. W. W. Norton and Company, Inc., New York: 1999.
Jones, Gareth Stedman. "Working-class culture and working-class politics in London, 1870-1900: Notes on the remaking of a working class" in Languages of class: studies in English working class history, 1832-1982. Cambridge University Press, Cambridge: 1983.
Smalley, Donald ed.
Poems of Robert Browning
. Houghton Mifflin Company Riverside Editions, Boston: 1956.
Swinburne, Algernon Charles. "Ave Atque Vale, In Memory of Charles Baudelaire." [text on U. of Toronto site].
Last modified 17 December 2003 |
|135–185 million Communities majorly in Iran,Afghanistan and also in Turkey, Iraq, Syria, Azerbaijan, Tajikistan, Uzbekistan, Turkmenistan, Russia, Kazakhstan, Georgia, Armenia, Oman, China (Xinjiang), India, Pakistan, United Kingdom, Germany and United States.|
|Regions with significant populations|
|Iran and Iranian Plateau, Anatolia, South Asia, Central Asia, the Caucasus and as immigrant communities in North America and Western Europe.|
|Related ethnic groups|
Other Indo-Iranian peoples
The Iranian peoples or Iranic peoples are an Indo-European ethno-linguistic group that comprise the speakers of Iranian languages, a major branch of the Indo-European language family, as such forming a branch of the Indo-European-speaking peoples. Their historical areas of settlement were on the Iranian plateau (mainly Iran) and certain neighbouring areas of Central Asia (such as Afghanistan, Tajikistan, Uzbekistan, Pakistan West of the River Indus, northern Iraq and eastern Turkey, and scattered part of the Caucasus Mountains) reflecting changing geopolitical range of the Persian empires and the Iranian history. Their current distribution spreads across the Iranian plateau, and stretches from Pakistan's Indus River in the east to eastern Turkey in the west, and from Central Asia and the Caucasus in the north to the Persian Gulf in the south – a region that is sometimes called the Iranian cultural continent, or Greater Persia by scholars, and represents the extent of the Iranian languages and influence of the Persian people, through the geopolitical reach of the Persian empire.
The Iranian group emerges from an earlier Iranian group during the Late Bronze Age, and it enters the historical record during the Early Iron Age.
The Iranians comprise the Persians, Medes, Scythians, Bactrians, Parthians, Sarmatians, Alans, Ossetians, Cimerians and their sub-groups. The Iranians had domesticated horses, had travelled far and wide, and from the late 2nd millennium BCE to early 1st millennium BCE they had migrated to, and settled on, the Iranian Plateau. They moved into the Zagros Mountains (inhabited by Gutians, Kassites and others, home of the Mannaean kingdom) above the indigenous non Iranian Elamite Kingdom. For approximately three centuries after arriving in the region, the Medes and Persians fell under the domination of the Assyrian Empire (911–609 BCE), based in nearby Mesopotamia. In 646 BCE, Susa and many other cities of Elam were plundered and wrecked by Ashurbanipal, King of Assyria, allowing the Iranian peoples to become the predominant group in Iran. After the death of Ashurbanipal in 627 BCE, the Assyrian Empire began to unravel due to a series of bitter civil wars. In 616 BCE the Median king Cyaxares threw off the Assyrian yoke, united the Medes and Persians, and in alliance with Nabopolassar of Babylon and the Scythians, attacked the civil war ridden Assyrian Empire. By 609 BCE, the Assyrians and their Egyptian allies had been defeated. This began the Iranian domination in the Iranian Plateau. Persians formed the Achaemenid Empire by the 6th century BCE, while the Scythians dominated the Eurasian steppe. With numerous artistic, scientific, architectural and philosophical achievements and numerous kingdoms and empires that bridged much of the civilized world in antiquity, the Iranian peoples were often in close contact with the Greeks, Romans, Egyptians, Indians, and Chinese. The various religions of the Iranian peoples, including Zoroastrianism, Mithraism and Manichaeism, are believed by some scholars to have been significant early philosophical influences on Christianity and Judaism.
The term Iranian is derived from the Old Iranian ethnical adjective Aryana which is itself a cognate of the Sanskrit word Arya. The name Iran is from Aryānām; lit: "[Land] of the Aryans". The old Proto-Indo-Iranian term Arya, per Thieme meaning "hospitable", is believed to have been one of the self-referential terms used by the Aryans, at least in the areas populated by Aryans who migrated south from Central Asia. Another meaning for Aryan is noble. In the late part of the Avesta (Vendidad 1), one of their homelands was referred to as Airyanem Vaejah. The homeland varied in its geographic range, the area around Herat (Pliny's view) and even the entire expanse of the Iranian plateau (Strabo's designation).
The academic usage of the term Iranian is distinct from the state of Iran and its various citizens (who are all Iranian by nationality and thus popularly referred to as Iranians) in the same way that Germanic people is distinct from Germans. Many citizens of Iran are not necessarily "Iranian people" by virtue of not being speakers of Iranian languages. Unlike the various terms connected with the Aryan arya- in Old Indian, the Old Iranian term has solely an ethnic meaning and there can be no doubt about the ethnic value of Old Iran. arya (Benveniste, 1969, I, pp. 369 f.; Szemerényi; Kellens).
The Avesta clearly uses airya as an ethnic name (Vd. 1; Yt. 13.143-44, etc.), where it appears in expressions such as airyāfi; daiŋˊhāvō "Iranian lands, peoples," airyō.šayanəm "land inhabited by Iranians," and airyanəm vaējō vaŋhuyāfi; dāityayāfi; "Iranian stretch of the good Dāityā," the river Oxus, the modern Āmū Daryā.
The term "Ariya" appears in the royal Old Persian inscriptions in three different contexts: 1) As the name of the language of the Old Persian version of the inscription of Darius the Great in Behistun; 2) as the ethnic background of Darius in inscriptions at Naqsh-e-Rostam and Susa (Dna, Dse) and Xerxes in the inscription from Persepolis (Xph) and 3) as the definition of the God of Iranian people, Ahuramazda, in the Elamite version of the Behistun inscription. For example in the Dna and Dse Darius and Xerxes describe themselves as "An Achaemenian, A Persian son of a Persian and an Aryan, of Aryan stock". Although Darius the Great called his language the Iranian language, modern scholars refer to it as Old Persian because it is the ancestor of modern Persian language.
The Old Persian and Avestan evidence is confirmed by the Greek sources". Herodotus in his Histories remarks about the Iranian Medes that: "These Medes were called anciently by all people Arians; " (7.62). In Armenian sources, the Parthians, Medes and Persians are collectively referred to as Iranians. Eudemus of Rhodes apud Damascius (Dubitationes et solutiones in Platonis Parmenidem 125 bis) refers to "the Magi and all those of Iranian (áreion) lineage"; Diodorus Siculus (1.94.2) considers Zoroaster (Zathraustēs) as one of the Arianoi.
The name of Ariana is further extended to a part of Persia and of Media, as also to the Bactrians and Sogdians on the north; for these speak approximately the same language, with but slight variations.
— Geography, 15.8
The trilingual inscription erected by Shapur's command gives a more clear description. The languages used are Parthian, Middle Persian and Greek. In Greek, the inscription says: "ego ... tou Arianon ethnous despotes eimi"("I am lord of the kingdom (Gk. nation) of the Aryans") which translates to "I am the king of the Iranian people". In the Middle Persian, Shapour states: "ērānšahr xwadāy hēm" and in Parthian he states: "aryānšahr xwadāy ahēm".
The Bactrian language (a Middle Iranian language) inscription of Kanishka the founder of the Kushan empire at Rabatak, which was discovered in 1993 in an unexcavated site in the Afghanistan province of Baghlan, clearly refers to this Eastern Iranian language as Arya. In the post-Islamic era, one can still see a clear usage of the term Iran in the work of the 10th-century historian Hamzeh Isfahani. In his book the history of Prophets and Kings writes: "Aryan which is also called Pars (Persia) is in the middle of these countries and these six countries surround it because the South East is in the hands China, the North of the Turks, the middle South is India, the middle North is Rome, and the South West and the North West is the Sudan and Berber lands". All this evidence shows that the name arya "Iranian" was a collective definition, denoting peoples (Geiger, pp. 167 f.; Schmitt, 1978, p. 31) who were aware of belonging to the one ethnic stock, speaking a common language, and having a religious tradition that centered on the cult of Ahura Mazdā.
History and settlement
The language referred to as Proto-Indo-European (PIE): is ancestral to Diba and the Celtic, Italic (including Romance), Germanic, Baltic, Slavic, Indo-Iranian, Albanian, Armernian, Greek, and Tocharian languages.
'There is an agreement that the PIE community split into two major groups from wherever its homeland was situated (its location is unknown), and whenever the timing of its dispersal (also unknown). One headed west for Europe and became speakers of Indo-European (all the languages of modern Europe save for Basque, Hungarian, and Finnish) while others headed east for Eurasia to become Indo-Iranians. The Indo-Iranians were a community that spoke a common language prior to their branching off into the Iranian and Indo-Aryan languages. Iranian refers to the languages of Iran (Iranian), Pakistan (Balochi and Pashto), Afghanistan (Pashto and Dari), and Tadjikistan (Tajiki) and Indo-Aryan, Sanskrit, Urdu and its many related languages.' – (Carl C. Lamberg-Karlovsky: Case of the Bronze Age)
By the early 1st millennium, Ancient Iranian peoples such as Medes, Persians, Bactrians, Parthians and Scythians populated the Iranian plateau, and other Scythian tribes, along with Cimmerians, Sarmatians and Alans populated the steppes north of the Black Sea. The Saka, Scythian, tribes spread as far west as the Balkans and as far east as Xinjiang. Scythians as well formed the Indo-Scythian Empire, and Bactrians formed a Greco-Bactrian Kingdom founded by Diodotus I, the satrap of Bactria. The Kushan Empire, with Bactrian roots/connections, once controlled much of Pakistan, some of Afghanistan and Tajikistan. The Kushan elite (who the Chinese called the Yuezhi) were either a Tocharian-speaking (another Indo-European branch) people or an Eastern Iranian language-speaking people.
The division into an "Eastern" and a "Western" group by the early 1st millennium is visible in Avestan vs. Old Persian, the two oldest known Iranian languages. The Old Avestan texts known as the Gathas are believed to have been composed by Zoroaster, the founder of Zoroastrianism, with the Yaz culture (c. 1500–1100 BCE) as a candidate for the development of Eastern Iranian culture.
Western Iranian peoples
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During the 1st centuries of the first millennium BCE, the ancient Persians established themselves in the western portion of the Iranian plateau and appear to have interacted considerably with the Elamites and Babylonians, while the Medes also entered in contact with the Assyrians. Remnants of the Median language and Old Persian show their common Proto-Iranian roots, emphasized in Strabo and Herodotus' description of their languages as very similar to the languages spoken by the Bactrians and Soghdians in the east. Following the establishment of the Achaemenid Empire, the Persian language (referred to as "Farsi" in Persian) spread from Pars or Fars Province to various regions of the Empire, with the modern dialects of Iran, Afghanistan (also known as Dari) and Central-Asia (known as Tajiki) descending from Old Persian.
Old Persian is attested in the Behistun Inscription (c. 519 BCE), recording a proclamation by Darius the Great. In southwestern Iran, the Achaemenid kings usually wrote their inscriptions in trilingual form (Elamite, Babylonian and Old Persian) while elsewhere other languages were used. The administrative languages were Elamite in the early period, and later Imperial Aramaic.
The early inhabitants of the Achaemenid Empire appear to have adopted the religion of Zoroastrianism. The Baloch who speak a west Iranian language relate an oral tradition regarding their migration from Aleppo, Syria around the year 1000 CE, whereas linguistic evidence links Balochi to Kurmanji, Soranî, Gorani and Zazaki.
Eastern Iranian peoples
While the Iranian tribes of the south are better known through their texts and modern counterparts, the tribes which remained largely in the vast Eurasian expanse are known through the references made to them by the ancient Greeks, Persians, Indo-Aryans as well as by archaeological finds. Many ancient Sanskrit texts make references to tribes like Sakas, Paradas, Kambojas, Bahlikas, Uttaramadras, Madras, Lohas, Parama Kambojas, Rishikas, Tukharas or Tusharas etc. and locate them in the (Uttarapatha) (north-west) division, in Central Asia, around Hindukush range in northern Pakistan. The Greek chronicler, Herodotus (5th century BCE) makes references to a nomadic people, the Scythians; he describes them as having dwelt in what is today southern Russia.
It is believed that these Scythians were conquered by their eastern cousins, the Sarmatians, who are mentioned by Strabo as the dominant tribe which controlled the southern Russian steppe in the 1st millennium CE. These Sarmatians were also known to the Romans, who conquered the western tribes in the Balkans and sent Sarmatian conscripts, as part of Roman legions, as far west as Roman Britain.
The Sarmatians of the east became the Alans, who also ventured far and wide, with a branch ending up in Western Europe and North Africa, as they accompanied the Germanic Vandals during their migrations. The modern Ossetians are believed to be the sole direct descendants of the Alans, as other remnants of the Alans disappeared following Germanic, Hunnic and ultimately Slavic migrations and invasions. Another group of Alans allied with Goths to defeat the Romans and ultimately settled in what is now called Catalonia (Goth-Alania).
Some of the Saka-Scythian tribes in Central Asia would later move further southeast and invade the Iranian plateau, large sections of present day Afghanistan and finally deep into present day Pakistan (see Indo-Scythians). Another Iranian tribe related to the Saka-Scythians were the Parni in Central Asia, and who later become indistinguishable from the Parthians, speakers of a northwest-Iranian language. Many Iranian tribes, including the Khwarazmians, Massagetae and Sogdians, were assimilated and/or displaced in Central Asia by the migrations of Turkic tribes emanating out of Xinjiang and Siberia.
The most dominant surviving Eastern Iranian peoples are represented by the Pashtuns, whose origins are generally believed to be from the Sulaiman Mountains, from which they began to spread until they reached as far west as Herat, north to areas of southern and eastern Afghanistan; and as eastward towards the Indus. The Pashto language shows affinities to the Avestan and Bactrian.
The modern Sarikoli in southern Xinjiang and the Ossetians of the Caucasus are remnants of the various Saka tribes. The modern Ossetians claim to be the descendants of the Alano-Sarmatians and their claims are supported by their Northeast Iranian language, while culturally the Ossetians resemble their Caucasian neighbors, the Kabardians and Circassians. Various extinct Iranian people existed in the eastern Caucasus, including the Azaris, while some Iranian people remain in the region, including the Talysh and the Tats (including the Judeo-Tats, who have relocated to Israel), found in Azerbaijan and as far north as the Russian republic of Dagestan. A remnant of the Sogdians is found in the Yaghnobi speaking population in parts of the Zeravshan valley in Tajikistan.
Later developments
Starting with the reign of Omar in 634 CE, Muslim Arabs began a conquest of the Iranian plateau. The Arabs conquered the Sassanid Empire of the Persians and seized much of the Byzantine Empire populated by the Kurds and others. Ultimately, the various Iranian people, including the Persians, Azaries, Kurds, Baluchis and Pashtuns, converted to Islam. The Iranian people would later split along sectarian lines as the Persians (and later the Hazara) adopted the Shi'a sect. As ancient tribes and identities changed, so did the Iranian people, many of whom assimilated foreign cultures and people.
Later, during the 2nd millennium CE, the Iranian people would play a prominent role during the age of Islamic expansion and empire. Saladin, a noted adversary of the Crusaders, was an ethnic Kurd, while various empires centered in Iran (including the Safavids) re-established a modern dialect of Persian as the official language spoken throughout much of what is today Iran and adjacent parts of Central Asia. Iranian influence spread to the Ottoman Empire, where Persian was often spoken at court, as well to the court of the Mughal Empire. All of the major Iranian people reasserted their use of Iranian languages following the decline of Arab rule, but would not begin to form modern national identities until the 19th and early 20th centuries (just as Germans and Italians were beginning to formulate national identities of their own).
The following either partially descend from Iranian people or are sometimes regarded as possible descendants of ancient Iranian people:
- Azeris: Although Azeris speak a Turkic language (modern Azerbaijani language), they are believed to be primarily descendants of ancient Iranians. Thus, due to their historical ties with various ancient Iranians, as well as their cultural ties to Persians, the Azeris are often associated with the Iranian people (see Origin of Azerbaijani people and the Iranian theory regarding the origin of the Azerbaijanis for more details).
- Uzbeks: The modern Uzbek people are believed to have both Iranian and Turkic ancestry. "Uzbek" and "Tajik" are modern designations given to the culturally homogeneous, sedentary population of Central Asia. The local ancestors of both groups – the Turkic-speaking Uzbeks and the Iranian-speaking Tajiks – were known as "Sarts" ("sedentary merchants") prior to the Russian conquest of Central Asia, while "Uzbek" or "Turk" were the names given to the nomadic and semi-nomadic populations of the area. Still today, modern Uzbeks and Tajiks are known as "Sarts" to their Turkic neighbours, the Kazakhs and the Kyrgyz. The ancient Soghdians and Bactrians are among their ancestors. Culturally, the Uzbeks are closer to their sedentary Iranian-speaking neighbours rather than to their nomadic and semi-nomadic Turkic neighbours. Some Uzbek scholars, i.e. Ahmadov and Askarov, favour the Iranian origin theory.
- The native name of Yakuts is Sakha, very similar to the Sakkas, proposing Yakuts to be related of descendants of Scythians, specifically Sakkahs.
- Volga TatarsMany are mixed from Volga bulgars. The reasons are same with Bulgarians, and the putative claim on the Iranian origin of bulgars.
- A few linguists suggest that the names of the South Slavic people, the Serbs and Croats are of Iranian origin. Those who entertain such a connection propose that the Sarmatian Serboi and Kharoti tribes might have migrated from the Eurasian steppe lands to eastern Europe, and assimilated with the numerically superior Slavs, passing on their name. Iranian-speaking people did inhabit parts of the Balkans in late classical times, and would have been encountered by the Slavs. However, direct linguistic, historical or archaeological proof for such a theory is lacking. (See also: Theories on the origin of Serbs and Theories on the origin of Croats)Ultimately, Montenegrins and Bosniaks may be counted to this theory.
- Some modern Bulgarian historians claim that the Bulgars were of Iranian origin and that they migrated to Europe from the region of today's northern Afghanistan – Hindukush mountains, from the Kingdom of Balhara. Their claims are based on medieval Armenian sources, the writings of ancient historians ("Ashharatsuyts" by Anania Shirakatsi; Agathias of Myrina, Theophylact Simocatta, Michael the Syrian) archaeological findings in modern Bulgaria, the similarities with Iranian languages (place names, people names, and Iranian words in modern Bulgarian), similarities with culture (e.g.: some buildings in Pliska were built in a Zoroastrian fashion; similarities in traditional music, dancing and carpet making) and the very close similarity of the DNA of Pamirian/Iranian people with that of modern Bulgarians After their arrival on Balkans, the Bulgars subjugated and then formed an alliance with the local Slavs and formed the Bulgarian nation. Ultimately, Slavic Macedonians could be counted due to their close linguistic affinities with the standard Bulgarian language.
|english||persian||zazaki||(Kurdish) Kurmanji / Sorani||bulgarian|
|i know||midânam||ez dizono||ez dizanim /min azanim||az znam|
|you know||midâni||ti dizana||tu dizanî / to azanit||ti znayş|
|i don't know||nemidânam||ez nizon||ez nizanim / min nazanim||az neznam|
|you don't know||nemidâni||ti nizona||tu nizanî / to nazanit||ti niznayş|
|a dog||sag||kûtik||kûtchik / sag||kutche (kûçe)|
- Indo-Aryan speakers
- Speakers of Indo-Aryan languages share linguistic affinities with speakers of Iranian languages, which suggests a degree of historical interaction between these two groups.
- Brahui people in Pakistan are speakers of a language classified as Dravidian, although culturally there is considerable Iranic influence among Brahui populations.
- Uralic speakers
- Many Volga Finns may be of part Iranian admixture due to Bulgar invasion of the Volga basin, if they (Bulgars) were Iranian people.
- Hungarians have long prided themselves as Scythians in the past, Scythians being an Iranian people, prior to the Finno-Ugric/Uralic theory. It's possible they've undergone a language shift. In a Magyar folkore suggests Iranian admixture among Hungarian, when Hunor and Magor marry princesses who were Alans, another Iranian people. Jassic people of Hungary are of Ossetian origin. The Szekely are possibly of Iranian origin, as their name is similar to Sakka.
- Shirazis:The Shirazi are a sub-group of the Swahili people living on the Swahili Coast of East Africa, especially on the islands of Zanzibar, Pemba and Comoros. Local traditions about their origin claim they are descended from merchant princes from Shiraz in Persia who settled along the Swahili Coast.
There are an estimated 150 to 200 million native speakers of Iranian languages, the five major groups of Persians, Lurs, Kurds, Baloch, and Pashtuns accounting for about 90% of this number. Currently, most of these Iranian peoples live in Iran, the Caucasus (mainly Ossetia, other parts of Georgia, and Azerbaijan), Iraqi Kurdistan and Kurdish majority populated areas of Turkey, Iran and Syria, Afghanistan, Tajikistan, and Uzbekistan.
The following is a list of peoples that speak Iranian languages with the respective groups's core areas of settlements and their estimated sizes (in millions):
|Persian-speaking peoples||Iran, Afghanistan, Tajikistan, Uzbekistan, Iraq, Bahrain||
|Kurds||Turkey, Iran, Iraqi Kurdistan, Syria||
|Baluchis||Pakistan, Iran, Afghanistan||
|Gilakis & Mazanderanis||Iran||
|Lurs & Bakhtiaris||Iran||
|Pamiri people||Tajikistan, Afghanistan, China (Xinjiang), Pakistan||
|Ossetians||South Ossetia, Georgia,
Russia (North Ossetia), Hungary
|Yaghnobi||Uzbekistan and Tajikistan (Zerafshan region)||
It is largely through linguistic similarities that the Iranian people have been linked, as many non-Iranian people have adopted Iranian languages and cultures. However, other common traits have been identified as well, for example, a stream of common historical events have often linked the southern Iranian people, including Hellenistic conquests, the various empires based in Persia, Arab Caliphates and Turkic invasions.
Like other Indo-Europeans, the early Iranians practiced ritual sacrifice, had a social hierarchy consisting of warriors, clerics and farmers and poetic hymns and sagas to recount their deeds.
Following the Iranian split from the Indo-Iranians, the Iranians developed an increasingly distinct culture. Various common traits can be discerned among the Iranian people. For example, the social event Norouz is an Iranian festival that is practiced by nearly all of the Iranian people as well as others in the region. Its origins are traced to Zoroastrianism and pre-historic times.
Some Iranian cultures exhibit traits that are unique unto themselves. The Pashtuns adhere to a code of honor and culture known as Pashtunwali, which has a similar counterpart among the Baloch, called Mayar, that is more hierarchical.
The early Iranian people worshipped various deities found throughout other cultures where Indo-European immigrants established themselves. The earliest major religion of the Iranian people was Zoroastrianism, which spread to nearly all of the Iranian people living in the Iranian plateau. Other religions that had their origins in the Iranian world were Mithraism, Manichaeism, and Mazdakism, among others.
Modern speakers of Iranian languages mainly follow Islam. Some follow Judaism, Christianity, Zoroastrianism, and the Bahá'í Faith, with an unknown number showing no religious affiliation. Overall the numbers of Sunni and Shia among the Iranian people are equally distributed. Most Kurds, Tajiks, Pashtuns, and Baluch are Sunni Muslims, while the remainder are mainly Twelver Shi'a, comprising mostly Persians in Iran, and Hazaras in Afghanistan. Zazas in Turkey are largely Alevi, while the Pamiri peoples in Tajikistan and China are nearly all Ismaili. The Christian community is mainly represented by the Armenian Apostolic Church, followed by the Russian Orthodox and Georgian Orthodox Ossetians followed by Nestorians. Judaism is followed mainly by Persian Jews, Kurdish Jews, Bukharian Jews (of Central Asia) and the Mountain Jews (of the Caucasus), most of whom are now found in Israel. The historical religion of the Persian Empire was Zoroastrianism and it still has a few thousand followers, mostly in Yazd and Kerman. They are known as the Parsis in the Indian subcontinent, where many of them fled in historic times following the Arab conquest of Persia, or Zoroastrians in Iran. Another ancient religion is the Yazidi faith, followed by some Kurds in northern Iraq, as well as the majority of the Kurds in Armenia.
Cultural assimilation
In matters relating to culture, the various Turkic-speaking ethnic groups of Iran (notably the Azerbaijani people) and Afghanistan (Uzbeks and Turkmen) are often conversant in Iranian languages, in addition to their own Turkic languages and also have Iranian culture to the extent that the term Turko-Iranian can be applied. The usage applies to various circumstances that involve historic interaction, intermarriage, cultural assimilation, bilingualism and cultural overlap or commonalities.
Notable among this synthesis of Turko-Iranian culture are the Azeris, whose culture, religion and significant periods of history are linked to the Persians. Certain theories and genetic tests suggest that the Azeris are genetically more Iranian than Turkic.
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R1 is more closely linked to Iranians, while R1b is linked to Europeans.
Haplogroup J2 especially the clade J2a is frequently found among almost all groups of Iranian people. In comparison with the haplogroup R1a1, J2 is not restricted to geographically eastern and western Iranian populations, but also found among north-western and south-western Iranian populations such as the Bakhtiaris and Mazanderani, as well as geographically north-western Iranian Ossetians. Despite its supposed origin in the fertile crescent, J2a is also found among Iranian populations in the east such as the Yagnobi which are of Soghdian origin as well as the Parsis of India. Beside the relatively high percentage among the Yagnobis in Central Asia, other Iranian populations tend to have a higher frequency of J2a when compared to neighboring Turkic populations. The relatively strong presence of J2a among Ossetians as well as Yagnobis proves distant from the supposed Mesopotamian origin region of J2, are carriers of this Haplogroup.
In the Indo-Iranian context, the occurrence of J2a in South Asia is limited to caste populations, with the highest frequencies found among northern areas of South Asia. Compared with R1a1, J2a shows a more conservative distribution, stronger limited to Indo-Iranian origin groups.
Many Haplotypes of Y-chromosomal Haplogroup R have been found throughout the Iranian Plateau, and it has been suggested that this Haplogroup may have had its origins in Iran. Cambridge University geneticist Toomas Kivisild has suggested : "Given the geographic spread and STR diversities of sister clades R1 and R2, the latter of which is restricted to India, Pakistan, Iran, and southern central Asia, it is possible that southern and western Asia were the source for R1 and R1a differentiation."(Kivisild et al. 2003). A similar conclusion was given by population geneticist Miguel Regueiro in the Journal of Human Heredity (Regueiro et al. Human Heredity vol. 61 (2006), pp. 132–143)
Genetic studies conducted by Cavalli-Sforza have revealed that Iranians have weak correlation with Near Eastern groups, and are closer to surrounding Indo-Europeans speaking populations. This study is partially supported by another one, based on Y-Chromosome haplogroups.
The findings of this study reveal many common genetic markers found among the Iranian people from the Tigris river of Iraq to the Indus of Pakistan. This correlates with the Iranian languages spoken from the Caucasus to Kurdish areas in the Zagros region and eastwards to western Pakistan and Tajikistan and parts of Uzbekistan in Central Asia. The extensive gene flow is perhaps an indication of the spread of Iranian-speaking people, whose languages are now spoken mainly on the Iranian plateau and adjacent regions.
Another recent study of the genetic landscape of Iran was done by a team of Cambridge geneticists led by Dr. Maziar Ashrafian Bonab (an Iranian Azarbaijani). Bonab remarked that his group had done extensive DNA testing on different language groups, including Indo-European and non Indo-European speakers, in Iran. The study found that the Azerbaijanis of Iran do not have a similar FSt and other genetic markers found in Anatolian and European Turks. However, the genetic Fst and other genetic traits like MRca and mtDNA of Iranian Azeris were identical to Persians in Iran. Azaris of Iran also show very close genetic ties to Kurds.
See also
Literature and further reading
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The name ἀντίχριστος signifies an opposer of Christ. It is used only by John in his first and second epistles, though those opposed to Christ are referred to by others under different names. It is important to distinguish between an antichrist and the antichrist. John says, "as ye have heard that antichrist shall come, even now are there many antichrists;" whereas "he is the antichrist that denieth the Father and the Son." 1 John 2:18, 22. He is the consummation of the many antichrists. To deny Jesus Christ come in the flesh is the spirit or power of the antichrist, but it eventuates in a departure from the special revelation of Christianity: 'they went out from us.' 1 John 2:19; 1 John 4:3; 2 John 7. Now this clears the ground at once of much that has obscured the subject. For instance, many have concluded that Popery is the antichrist, and have searched no farther into the question, whereas the above passage refutes this conclusion, for Popery does not deny the Father and the Son; and, in Revelation 17, 18, Popery is pointed out as quite distinct from 'the false prophet,' which is another name for the antichrist. It is fully granted that Popery is anti-christian, and a Christ-dishonouring and soul-deceiving system; but where God has made a distinction we must also do so. Besides Popery there were and there are many antichrists, which, whatever their pretensions, are the enemies of Christ, opposers of the truth, and deceivers of man.
As to the Antichrist, it should be noticed that John makes another distinction between this one and the many. He speaks of the many as being already there, whereas the one was to come; and if we turn to 2 Thess. 2:3-12 we read of something or some one that hinders that wicked or lawless one being revealed, although the mystery of iniquity was already at work. Now there has been no change of dispensation since this epistle was written, and John wrote much later, from which we learn that the revelation of the antichrist is still future, though doubtless the mystery of iniquity is getting ripe for his appearing; that which hindered and still hinders the manifestation of the antichrist is doubtless the presence of the Holy Spirit on earth. He will leave the earth at the rapture of the saints.
This passage in Thessalonians gives us further particulars as to this MAN OF SIN. His coming is after the working of Satan, that is, he will be a confederate of Satan, and be able to work signs and lying wonders with all deceit of unrighteousness in them that perish. Those that have refused the truth will then receive the lie of this wicked one. We get further particulars in Rev. 13:11-18, where the anti-christian power or kingdom is described as a beast rising out of the earth, having two horns as a lamb, but speaking as a dragon. Here again we read that he will do great wonders, making fire come down from heaven, with other signs or miracles.
In the description in Thessalonians he opposeth himself against all that is called God or that is worshipped, and sits down in the temple of God, and sets forth himself as God. The Jews will receive him as their Messiah, as we read in John 5:43. In the above passage in the Revelation this counterfeit of Christ's kingdom is openly idolatrous. He directs the dwellers on the earth to make an image of the beast (named in ver. 1, the future head of the resuscitated Roman empire) to which image he gives breath, that it should speak, and persecutes those who will not worship the image. He also causes all to receive a mark on their hand or their forehead that they may be known to be his followers; and that none else should be able to buy or sell. We thus see that in the Revelation the anti-christian power called also 'the false prophet' will work with the political head, and with Satan — a trinity of evil — not only in deceiving mankind, but also, in Rev. 16:13-16, gathering together by their influence the kings of the earth to the battle of that great day of God Almighty. The three are cast into the lake of fire Rev. 19:20; Rev. 20:10.
In the O.T. we get still another character of this wicked one. In Dan. 11:36-39 he is called 'king.' Here he exalts himself and speaks marvellous things against the God of gods. He will not regard the God of his fathers (pointing out that he will be a descendant of Israel, probably from the tribe of Dan, cf. Gen. 49:17), nor "the desire of women" (i.e. the Messiah, of whom every Jewess hoped to be the mother): he exalts himself above all. Here again he is an idolater, honouring a god that his fathers knew not. In Zech. 11:15-17 he is referred to as the foolish and idol shepherd, who cares not for the flock, in opposition to the Lord Jesus the good Shepherd.
This man of sin will 'do according to his own will' — just what the natural man ever seeks to do. In contrast to this the blessed Lord was obedient, and came not to do His own will. May His saints be ever on the watch against the many false prophets in the world, 1 John 4:1, and be loyal to their absent Lord, behold His beauty in the sanctuary, and reproduce Him more down here in their earthen vessels.
Strictly, those opposed to the inculcation of good works from a perverted view of the doctrines of grace; but the term is also falsely applied to those who know themselves free through the death of Christ from the law as given by Moses. Rom. 7:4; Gal. 2:19. One has but to read carefully the epistle to the Galatians to see that for Gentile believers to place themselves under the law is to fall from grace; and Paul exhorted them to be as he was, for he was (though a Jew by birth) as free from the law by the death of Christ as they were as Gentiles. They had not injured him at all by saying he was not a strict Jew, Gal. 4:12: in other words, they may have called him an antinomian, as others have been called, whose walk has been the most consistent. To go back to the law supposes that man has power to keep it. For a godly walk the Christian must walk in the Spirit, and grace teaches that, "denying ungodliness and worldly lusts, we should live soberly, righteously, and godly in this present world." Titus 2:12. On the other hand, there have been, and doubtless are, some who deny good works as a necessary fruit of grace in the heart: grace, as well as everything else, has been abused by man. See LAW.
Antioch in Pisidia. [An'tioch in Pisid'ia]
A Roman colony of Phrygia in Asia Minor, founded by Seleueus Nicator. Its ruins are now called Yalobatch or Yalowaj. Paul's labour here was so successful that it roused the opposition of the Jews and he was driven to Iconium and Lystra; but he returned with Silas. Acts 13:14; Acts 14:19-21; 2 Tim. 3:11.
Antioch in Syria. [An'tioch in Syria]
This is memorable in the annals of the church as the city where the disciples were first called Christians, Where an assembly of Gentiles was gathered, and from which Paul and his companions went forth on their missionary journeys, and to which they twice returned. It formed a centre for their labours among the Gentiles, outside the Jewish influence which prevailed at Jerusalem; yet the church in this city maintained its fellowship with the assembly at Jerusalem and elsewhere. Acts 6:5; Acts 11:19-30; Acts 13:1; Acts 14:26; Acts 15:22-35; Acts 18:22; Gal. 2:11.
Antioch was once a flourishing and populous city, the capital of Northern Syria, founded by Seleueus Nicator, B.C. 300, in honour of his father Antiochus. It was afterwards adorned by Roman emperors, and was esteemed the third city was eventually the seat of the Roman proconsul of Syria. It stood on a beautiful spot on the river Orontes, where it breaks through between the mountains Taurus and Lebanon. It is now called Antakia 36 12', 36 10' E. It has suffered from wars and earthquakes, and is now a miserable place. Comparatively few antiquities of the ancient city are to be found, but parts of its wall appear on the crags of Mount Silpius.
There were several kings bearing this name who ruled over Syria, and though they are not mentioned by name in scripture, some of their actions are specified. These are so clear and definite that sceptics have foolishly said that at least this part of the prophecy of Daniel must have been written after the events! The Greek kingdom, the third of the four great empires, was, on the death of Alexander the Great, divided among his four generals, and this resulted principally in a series of kings who ruled in Egypt bearing the general name of PTOLEMY, and are called in scripture 'Kings of the South;' and another series, called 'Kings of the North,' who bore the general name of either SELEUCUS or ANTIOCHUS. Both the Ptolemies and the Seleucidae began eras of their own, and some of the kings of each era had to do with Palestine and the Jews. The following is a list of the kings, with the dates when they began to reign, noticing the principal events that were prophesied of them in Daniel 11.
320 Ptolemy I, Soter. He takes Jerusalem. Era of the Ptolemies begins.
312 SELEUCUS I, Nicator. He re-takes Palestine. Era of the Seleucidae begins.
283 Ptolemy II, Philadelphus. The O.T. translated into Greek.
280 ANTIOCHUS I, Soter.
261 ANTIOCHUS II, Theos. He was at war with Ptolemy, but peace was restored on condition that Antiochus should put away his wife Laodice and marry Berenice the daughter of Ptolemy. This was done, but on the death of Philadelphus he restored Laodice; but she, fearing another divorce, poisoned her husband, and then caused the death of Berenice and her son. See Dan. 11:6.
247 Ptolemy III, Euergetes. He revenged his sister's death, being 'a branch of her roots;' and carried off 40,000 talents of silver, etc. 'Shall enter into the fortress of the king of the north,' and carry away their precious vessels of silver and gold. Dan. 11:7-9.
246 SELEUCUS II, Callinicus.
226 SELEUCUS III, Ceraunus.
223 ANTIOCHUS III, the Great.
222 Ptolemy IV, Philopater. War between Ptolemy and Antiochus. Ptolemy recovers Palestine. Dan. 11:10-12.
205 Ptolemy V, Epiphanes (5 years old). Antiochus seized the opportunity of the minority of the king to regain the country. Dan. 11:16. He also joined with Philip of Macedonia to capture other portions of the dominions of Ptolemy. But Rome was now growing in power, and on being appealed to by Egypt for protection, Antiochus was told he must let Egypt alone. In the meantime an army from Egypt had re-taken Palestine; but Antiochus, on his return, again obtained the mastery there. Wishing to extend his dominions in the west he proposed that Ptolemy should marry his daughter Cleopatra, that she might serve her father's ends; but she was faithful to her husband. Daniel thus speaks of it: "He shall give him the daughter of women, corrupting her, but she shall not stand on his side, neither be for him." Dan. 11:17. Antiochus took many maritime towns, but after many encounters he was compelled by Rome to quit all Asia on that side of Mount Taurus, give up his elephants and ships of war and pay a heavy fine. Antiochus had great difficulty in raising the money, and on attempting to rob a temple at Elymais he was killed. Dan. 11:18, 19.
187 SELEUCUS IV, Philopator, succeeded. His principal work was the raising of money to pay the war-tax to Rome. He ordered Heliodorus to plunder the temple; but Heliodorus poisoned him. He was thus 'a raiser of taxes,' and was 'destroyed neither in anger, nor in battle.' Dan. 11:20. Heliodorus seized the crown but was destroyed by Antiochus IV.
181 Ptolemy VI, Philometor. He was a minor, under his mother and tutors.
175 ANTIOCHUS IV, Epiphanes. He was not the rightful heir. He 'obtained the kingdom by flatteries.' He called himself Epiphanes, which is 'illustrious;' but he was such 'a vile person' that people called him Epimanes, 'madman.' Dan. 11:21-24. He invaded Egypt and was at first successful: cf. Dan. 11:25, 26. The two kings entered into negotiations, though neither of them was sincere in what they agreed to: their hearts were to do mischief, and they 'tell lies at one table.' Dan. 11:27. Then Antiochus returned to his land with great riches: his heart was 'against the holy covenant,' and he entered Jerusalem and even into the sanctuary and took away the golden altar, the candlestick, the table of showbread, the censers of gold, and the other holy vessels and departed. 'At the appointed time he shall return and come toward the South,' Dan. 11:29; but he was stopped by Rome; 'ships of Chittim,' ships from Macedonia, came against him; and in great anger he returned and vented his wrath on Jerusalem.
He sent an army there with orders to slay all the men and sell the women and children for slaves. This was to a certain extent carried out. The walls were also thrown down and the city pillaged and then set on fire. He then decreed that the Jews should forsake their religion, and all should worship the heathen gods. To ensure this at Jerusalem with the few that still clung to the place, an image of Jupiter Olympius was erected in the temple and on an altar sacrifices were offered to this god. This was in B.C. 168 on the 25th of the month Chisleu. Daniel relates "They shall pollute the sanctuary of strength, and shall take away the daily sacrifice, and they shall place the abomination that maketh desolate." Dan. 11:31: cf. also Dan. 8:9-12 where the 'little horn' refers to Antiochus Epiphanes.
Bleek, Delitzsch, and others consider that in Dan. 8:14, the 2,300 'evening, morning,' margin, refer to the daily sacrifice, which is spoken of in Dan. 8:11, 12, 13; and that by 2,300 is meant 1,150 days: cf. also Dan. 8:26. The dedication of the temple was on the 25th of Chisleu, B.C. 165, and the desecration began some time in the year 168.
Dan. 11:32b, 33-35 refer to the change that soon took place under Judas Maccabeus and his brothers, commencing B.C. 166, and in 165 the temple was re-dedicated. In B.C. 164 ANTIOCHUS V. Eupator succeeded to the throne; and in 162 DEMETRIUS SOTER; but they were not powerful against Judaea, and in B.C. 161 an alliance was made by Judaea with Rome. The historical notices in Daniel end at Dan. 11:35.
It will be seen by the above that the records of history agree perfectly with the prophecy, as faith would expect them to do. It is only unbelief that has any difficulty in God foretelling future events. Without doubt some of the acts of Antiochus Epiphanes are types of the deeds of the future king of the North — referred to in other prophecies as 'the Assyrian' — in respect to the Jews and Jerusalem.
1. A Christian of Pergamos, who was martyred. Rev. 2:13.
2. Son of Herod the Great, but not called Antipas in the N.T. See HEROD.
The town to which Paul was taken in the night from Jerusalem on his way to Caesarea. Acts 23:31. It was built by Herod the Great in a well-watered spot surrounded by a wood, and named after his father. At Ras el-Ain, 32 6' N, 34 56' E, are ruins which are held to mark the spot. This is 5 or 6 miles nearer Jerusalem than Kefr Saba, which some associate with Antipatris, because Josephus says it was called Kapharsaba before its name was altered by Herod. The former place being nearer to Jerusalem removes the difficulty that some have felt as to the distance of Antipatris being too far to reach in a night ; this reduces it to about 36 miles, and it would be even less by cross roads.
The word antitype does not occur in the A.V., but the Greek word ἀντίτυπον occurs in Heb. 9:24, translated 'figures,' and in 1 Peter 3:21, translated 'like figure.' It is that which answers to a type, as a wax impression answers to a seal: if the device is sunk, the impression will be raised, or vice versa. To take a simple but beautiful example, a lamb was offered up for a burnt offering both morning and evening under the law; and in the N.T. we read, "Behold the Lamb of God, which taketh away the sin of the world." It is plain that the morning and evening lamb in Israel were types and the death of the Lord Jesus was the antitype. In Heb. 9:23, the 'heavenly things' are the type, and 'holy places,' Heb. 9:24, the antitype, or what corresponded to the pattern. In 1 Peter 3:21, eight souls were saved through water, of which baptism is the figure, or what answers to it. Doubtless there are many other antitypes in the N.T., but every antitype must have a type to which it corresponds, though the correspondence may not lie on its surface. Where scripture is silent as to types and antitypes the teaching of the Holy Spirit is needed, or grievous error may result in associating two things together which have no spiritual connection, though names and words may seem to correspond.
A tower or fortress built by Herod the Great near the temple at Jerusalem in which he placed a guard to watch over the approaches to the sacred edifice. Josephus (Wars v. 5, 8) says it was situated "at the corner of two cloisters of the court of the temple; of that on the west, and that on the north; it was erected upon a rock fifty cubits in height and was on a great precipice." Where this precipice was is not known, for it is a much disputed question upon what part of the temple area the temple was built. There is a tower, now called Antonia, on the N.W. angle, and there are indications of a similar one having stood on the S.E. angle.
A descendant of Benjamin. 1 Chr. 8:24.
Son of Coz, of the posterity of Judah. 1 Chr. 4:8.
The ape is not indigenous to Palestine; they were brought in the days of Solomon, with gold, silver, ivory and peacocks by the ships of Tarshish. The word goph may signify any of the monkey tribe. 1 Kings 10:22; 2 Chr. 9:21.
A Christian of Rome saluted by Paul as 'approved in Christ.' Rom. 16:10.
Apharsachites, Apharsathchites. [Aphar'sachites, Aphar'sathchites]
Some unknown Assyrian tribe sent as colonists to Samaria under Asnapper. Ezra 4:9; Ezra 5:6; Ezra 6:6.
An unknown Assyrian tribe as the preceding. Ezra 4:9.
1. Royal city of the Canaanites, the king of which was killed by Joshua, Joshua 12:18: probably the same as APHEKAH in Joshua 15:53. Not identified.
2. City in the north border of Asher, from which in the time of Joshua the inhabitants were not expelled. Joshua 13:4; Joshua 19:30: called APHIK in Judges 1:31. Identified with Afka at the foot of the Lebanon between Baalbek and Byblus.
3. Place where the Philistines encamped when Israel was defeated. 1 Sam. 4:1.
4. Where the Philistines encamped when Saul and Jonathan were killed. 1 Sam. 29:1. Perhaps the same as No. 3.
5. City, the wall of which falling killed 27,000 of the Syrians, 1 Kings 20:26, 30; 2 Kings 13:17. It is identified with Fik, 32 47' N, 35 41' E, on the great road between Damascus and Jerusalem.
A 'mighty man of power,' an ancestor of Saul. 1 Sam. 9:1.
The margin of Micah 1:10 explains the name as 'house of dust,' so that there is a play upon the word 'dust:' 'in the house of dust roll thyself in the dust.' The LXX read 'the house in derision.' It may refer to OPHRAH in Joshua 18:23; 1 Sam. 13:17, a city in the tribe of Benjamin.
Head of the eighteenth course of priests for service in the temple. 1 Chr. 24:15.
Another name for the REVELATION, q.v., being its Greek title ἀποκάλυψις.
The name given to those Books which were attached to the MSS copies of the LXX, but which do not form a part of the canon of scripture. The term itself signifies, 'hidden,' 'secret,' 'occult;' and, as to any pretence of being a part of scripture, they must be described as 'spurious.' There are such writings connected with both the Old and the New Testament, but generally speaking the term 'Apocrypha' refers to the O.T. (for those connected with the N. Test. see APOSTOLIC FATHERS. The O.T. books are:
1 I. Esdras.
2 II. Esdras.
5 Chapters of Esther, not found in the Hebrew nor Chaldee.
6 Wisdom of Solomon.
7 Jesus, son of Sirach; or Ecclesiasticus; quoted Ecclus.
8 Baruch, including the Epistle of Jeremiah.
9 Song of the Three Holy Children
10 The History of Susanna.
11 Bel and the Dragon.
12 Prayer of Manasseh.
13 I. Maccabees.
14 II. Maccabees.
The Council of Trent in A.D. 1546, professing to be guided by the Holy Spirit, declared the Apocrypha to be a part of the Holy Scripture. The above fourteen books formed part of the English Authorised Version of 1611, but are now seldom attached to the canonical books. Besides the above there are a few others, as the III., IV., and V. Maccabees, book of Enoch, etc., not regarded by any one as a part of scripture. It may be noticed
1. That the canonical books of the O.T. were written in Hebrew (except parts of Ezra and Daniel which were in Chaldee); whereas the Apocrypha has reached us only in Greek or Latin, though Jerome says some of it had been seen in Hebrew.
2. Though the Apocrypha is supposed to have been written not later than B.C. 30, the Lord never in any way alludes to any part of it; nor do any of the writers of the N.T., though both the Lord and the apostles constantly quote the canonical books.
3. The Jews did not receive the Apocrypha as any part of scripture, and to 'them were committed the oracles of God.'
4. As some of the spurious books were added to the LXX Version (the O.T. in the Greek) and to the Latin translation of the LXX, some of the early Christian writers were in doubt as to whether they should be received or not, and this uncertainty existed more or less until the before mentioned Council of Trent decided that the greater part of the Apocrypha was to be regarded as canonical. Happily at that time the Reformation had opened the eyes of many Christians to the extreme corruption of the church of Rome, and in rejecting the claims of that church they were also freed from its judgement as to the Apocryphal books.
5. The internal evidences of the human authorship of the Apocrypha ought to convince any Christian that it can form no part of holy scripture.
Expressions of the writers themselves show that they had no thought of their books being taken for scripture. There are also contradictions in them such as are common to human productions. Evil doctrines also are found therein: let one suffice: "Alms doth deliver from death, and shall purge away all sin." Tobit 12:9. The value of holy scripture as the fountain of truth is such that anything that might in any way contaminate that spring should be refused with decision and scorn. Some parts of the Apocryphal books may be true as history, but in every other respect they should be refused as spurious. Nor can it be granted that we need the judgement of the church, could a universal judgement be arrived at, as to what is to be regarded as the canon of scripture. The Bible carries its own credentials to the hearts and consciences of the saints who are willing to let its power be felt.
City of Macedonia, in the district of Mygdonia, some 28 miles from Amphipolis and 35 from Thessalonica, through which Paul and Silas passed. Acts 17:1.
A convert from Alexandria, an eloquent man and mighty in the scriptures, who, when only knowing the baptism of John, taught diligently the things of Jesus. At Ephesus he was taught more perfectly by Priscilla and Aquila. He laboured at Corinth, following the apostle Paul, who could hence say 'I have planted, Apollos watered,' and subsequently he greatly desired Apollos to revisit Corinth. His name is associated with that of Paul in connection with the party spirit at Corinth, which the apostle strongly rebuked; but from his saying he had 'transferred these things to himself and to Apollos,' it would appear that the Corinthians had local leaders, under whom they ranged themselves, whom he does not name; and that he taught them the needed lesson, and established the general principle by the use of his own name and that of Apollos rather than the names of their leaders. Acts 18:24; Acts 19:1; 1 Cor. 1:12; 1 Cor. 3:4-22; 1 Cor. 4:6; 1 Cor. 16:12; Titus 3:13.
The Greek translation of the Hebrew name ABADDON, which signifies 'destroyer.' He is king of the locusts of the bottomless pit, and ruler over the destroying agents that proceed from thence: it is one of the characters of Satan. Rev. 9:11.
Though the word 'apostasy' does not occur in the A.V., the Greek word occurs from which the English word is derived. In Acts 21:21 Paul was told that he was accused of teaching the Jews who were among the Gentiles to apostatise from Moses. Paul taught freedom from the law by the death of the Christ and this would appear to a strict Jew as apostasy. The same word is used in 2 Thess. 2:3, where it is taught that the day of the Lord could not come until there came 'the apostasy,' or the falling from Christianity in connection with the manifestation of the man of sin. See ANTICHRIST.
Though the general apostasy there spoken of cannot come till after the saints are taken to heaven, yet there may be, as there has been, individual falling away. See, for instance, Heb. 3:12; Heb. 10:26, 28, and the epistle of Jude. There are solemn warnings also that show that such apostasy will be more and more general as the close of the present dispensation approaches. 1 Tim. 4:1-3. Now a falling away necessarily implies a position which can be fallen from, a profession has been made which has been deliberately given up. This is, as scripture says, like the dog returning to his vomit, and the sow to her wallowing in the mire. It is not a Christian falling into some sin, from which grace can recover him; but a definite relinquishing of Christianity. Scripture holds out no hope in a case of deliberate apostasy, though nothing is too hard for the Lord.
The Greek word ἀπόστολος signifies 'a messenger,' 'one sent,' and is used in this sense for any messenger in 2 Cor. 8:23; Phil. 2:25; and as 'one sent' in John 13:16. It is also used in a much higher and more emphatic sense, implying a divine commission in the one sent, first of the Lord Himself and then of the twelve disciples whom He chose to be with Him during the time of His ministry here. The Lord in His prayer in John 17:18 said, "As thou hast sent me into the world, even so have I also sent them into the world." He was the Sent One, and in Heb. 3:1 it is written "Wherefore, holy brethren, partakers of the heavenly calling, consider the Apostle and High Priest of our profession, Jesus."* They were to consider this One who had been faithful, and who was superior to Moses, to the Aaronic priests, and to angels, and was in the glory. The ordering of a dispensation depended on the apostolic office as divinely appointed.
* The word 'Christ' is omitted by the Editors.
APOSTLES, THE TWELVE. The Lord appointed these "that they should be with him, and that he might send them forth to preach, and to have power to heal sicknesses, and to cast out demons," and also to carry out the various commissions given by Christ on earth. It will be seen by the lists that follow that Lebbaeus, Thaddaeus and Judas are the same person; and that Simon the Canaanite (Cananaean) and Simon Zelotes are the same; Peter is also called Simon; and Matthew is called Levi.
1 Peter and
3 James and
5 Philip and
7 Thomas and
10 and Lebbaeus.
11 Simon the Cananean and
12 Judas Iscariot.
12 Judas I.
11 Simon Zelotes.
12 Judas I.
11 Simon Z.
Peter is always named first; he with James and John was with the Lord on the mount of transfiguration and also with the Lord at other times, though no one apostle had authority over the others: they were all brethren and the Lord was their Master. Judas Iscariot is always named last. In Matthew the word 'and' divides the twelve into pairs, perhaps corresponding to their being sent out two and two to preach. Bartholomew and Simon Zelotes are not mentioned after their appointment except in Acts 1.
When the Lord sent the twelve out to preach He bade them take nothing with them, for the workman was worthy of his food: and on their return they confessed that they had lacked nothing. Their mission was with authority as the sent ones of the Lord; sicknesses were healed and demons cast out; and if any city refused to receive them it should be more tolerable for Sodom and Gomorrha in the day of judgement than for that city. Matt. 10:5-15.
They received a new mission from the Lord as risen: see Luke 24; John 20. And before the ascension the apostles were bidden to tarry at Jerusalem until they were endued with power from on high. This was bestowed at the descent of the Holy Spirit on the day of Pentecost. They are also viewed first among the gifts with which the church was endowed by the Head of the body when He ascended up on high. Eph. 4:8-11. These gifts were for "the perfecting of the saints, for the work of the ministry, for the edifying of the body of Christ." The mystery hitherto hid in God was now revealed to His holy apostles and prophets by the Spirit, namely, that the Gentiles should be joint heirs, and a joint body, and partakers of His promise in Christ Jesus. Eph. 3. Paul was the special vessel to make known this grace. His apostleship occupies a peculiar place, he having been called by the Lord from heaven, and being charged with the gospel of the glory. See PAUL.
On the death of Judas Iscariot, Matthias, an early disciple, was chosen in his place, for there must be (irrespective of Paul, who, as we have seen, held a unique place) twelve apostles as witnesses of His resurrection, Acts 1:22; Rev. 21:14 as there must still be twelve tribes of Israel. James 1:1 ; Rev. 21:12. At the conference of the church in Jerusalem respecting the Gentiles 'the apostles' took a prominent part, with the elders. Acts 15. How many apostles remained at Jerusalem is not recorded: we do not read of 'the twelve' after Acts 6. Tradition gives the various places where they laboured, which may be found under each of their names. Scripture is silent on the subject, in order that the new order of things committed to Paul might become prominent, as the older things connected with Judaism vanished away: cf. 2 Peter 3:15, 16.
There were no successors to the apostles: to be apostles they must have 'seen the Lord.' Acts 1:21, 22; 1 Cor. 9:1; Rev. 2:2. The foundation of the church was laid, and apostolic work being complete the apostles passed away, there remain however, in the goodness of God, such gifts as are needed "till we all come in the unity of the faith, and of the knowledge of the Son of God, unto a perfect man, unto the measure of the stature of the fulness of Christ." Eph. 4:12, 13.
This designation is applied to the early Christian writers, who had known the apostles, or had known those who had been acquainted with them.
1. BARNABAS; 2. CLEMENT; 3. HERMAS; are supposed to be the persons so named in the N.T.: see under their respective names.
4. POLYCARP, Bishop of Smyrna. He wrote an epistle to the Philippians about A.D. 125, Irenaeus says Polycarp was "instructed by the apostles, and was brought into contact with many who had seen Christ." He died a martyr's death. An ancient letter gives a particular account of his martyrdom.
5. IGNATIUS, Bishop of Antioch. Seven epistles are supposed to have been written by him, but they have been grossly interpolated; eight or nine others are wholly spurious. He was a martyr.
6. PAPIAS, Bishop of Hierapolis in Phrygia. He is said to have heard the apostle John. Various writings are attributed to him, but of which only fragments remain. He also died a martyr.
7. An unknown author of an eloquent and interesting epistle to Diognetus. Nearly all the above writings are very different from the scripture except where that is quoted. There is a deep dark line of demarcation between them and the writings which are inspired. Some of them however are found at the end of some of the Greek Testaments and were formerly read in the churches. Happily all these are now eliminated from any association with the N.T. Besides the above there are six apocryphal 'Gospels,' a dozen 'Acts,' four 'Revelations,' the 'Passing away of Mary,' etc.
This term is not used in scripture in the modern sense of a compounder of drugs for medicine; but in that of a compounder of ointments, etc., such as would now be called a 'perfumer,' as it is rendered in the margin of Ex. 30:25, where the holy anointing oil is an ointment compounded "after the art of the apothecary." The same was said of the holy incense. Ex. 30:35; Ex. 37:29. Asa was buried in a tomb filled with sweet odours and spices prepared by the apothecaries' art. 2 Chr. 16:14: cf. also Neh. 3:8. Spices were also carried to the tomb of the Lord to embalm His body.
Son of Nadab, of the tribe of Judah. 1 Chr. 2:30, 31.
It would appear from the arrangements made by Moses that some of the judges were accounted as judges of appeal, but that Moses himself, as having the mind of God, was the ultimate judge. Ex. 18:13-26. It is not probable, when the kingdom was established, that all causes were tried at Jerusalem; but only cases of appeal from the tribal judges; and it was such that Absalom alludes to in 2 Sam. 15:2, 3: see also Deut. 16:18. It is evident from Deut. 17:8-12 that the mind of God was to be sought where He put His name, if the matter was too hard for the judges. The Jewish writers say that before and after the time of Christ on earth, appeals could be carried through the various courts to the Grand Sanhedrim at Jerusalem.
In the case of Paul appealing to Caesar, it was not an appeal from a judgement already given, as is the case in what is now called an appeal; but Paul, knowing the deadly enmity of the Jews, and the corruption of the governors, elected to be judged at the court of Caesar, which, as a Roman, he had the right to do. Acts 25:11. There is One who "cometh to judge the earth: with righteousness shall he judge the world, and the people with equity." Ps. 98:9.
Appearing of Christ.
This is to be distinguished from Christ coming for His saints, though intimately connected with it, for He will bring them with Him. "When Christ, who is our life, shall appear, then shall ye also appear with him in glory." Col. 3:4. Here it is the manifestation of Christ with His own, to be followed by the setting up of His kingdom and the apportionment of rewards to His saints. 2 Cor. 5:10. The Lord's servant is exhorted by His appearing and His kingdom to preach the word, etc. 2 Tim. 4:1, 2. The saints will be associated with Christ in His judgements at His appearing. Jude 14, 15. Christ will execute judgement on the Beast and the False Prophet and the western powers. Also on the Assyrian and the eastern powers that will oppress the Jews. The Jews and the ten tribes will be restored to their land in blessing, ushering in the Millennium. See ADVENT, SECOND.
Probably the wife of Philemon, whom Paul addresses in that epistle, ver. 2.
Appii Forum. [Ap'pii For'um]
Station on the Appian Way, the main road from Rome to the Bay of Naples, where brethren went to meet Paul though 43 miles from Rome. Acts 28:15. The road was 18 to 22 feet wide, and parts of the ancient paving stones may still be seen. It was constructed by Appius Claudius, hence its name.
Apple, Apple Tree.
This is generally supposed to refer to the citron but apples grow in Palestine, and the Arabic name for the apple (tuffuh) differs little from the Hebrew word, tappuach. Others believe the quince is alluded to, which is fragrant and of a golden colour. Cant. 2:3, 5; Cant. 7:8; Cant. 8:5; Joel 1:12. In Prov. 25:11 "a word fitly spoken" is like some elegant device, as "apples of gold in pictures [or baskets] of silver."
Apple of the Eye.
1. ishon. Gesenius says this word signifies 'little man' and then 'the little man of the eye; 'that is, "the pupil of the eye in which, as in a mirror, a person sees his own image reflected in miniature." He says "this pleasing image is found in several languages." It is the part of the eye specially to be guarded: God preserved His own as the apple of His eye. Deut. 32:10; Ps. 17:8. His law should be kept as a precious thing. Prov. 7:2.
2. babah, the black or pupil of the eye, or, as others, 'the gate of the eye.' To touch God's people is touching the apple of His eye. Zech. 2:8.
3. bath, daughter. The sense is, Let not the apple (the daughter) of thine eye cease to shed tears. Lam. 2:18. In all places 'the apple of the eye' is a beautifully figurative expression for that which must be tenderly cherished as a most choice treasure.
The word chagorah signifies 'anything girded on.' When Adam and Eve had sinned they discovered that they were naked, and sewed fig-leaves together and made aprons, Gen. 3:7; but were soon conscious that this did not cover their nakedness, for when God called to them they owned that they were naked, and hid behind the trees. This teaches that nothing that man can devise can cover him from the eye of God. God clothed Adam and Eve with coats of skins; it was through death, typical of Christ Himself. In Acts 19:12 the word is σιμικίνθιον, and occurs but that once; it signifies a narrow apron or linen covering.
A converted Jew of Pontus, husband of Priscilla, whom Paul first met at Corinth. Acts 18:2. He and Paul worked together as tent-makers. Aquila and Priscilla had been driven from Rome as Jews by an edict of the emperor Claudius. They travelled with Paul to Ephesus, where they were able to help Apollos spiritually. Acts 18:18-26. They were still at Ephesus when Paul wrote 1 Corinthians (1 Cor. 16:19); and were at Rome when the epistle to the saints there was written, in which Paul said they had laid down their necks for his life, and that to them all the churches, with Paul, gave thanks. Rom. 16:3, 4. In Paul's last epistle he still sends his greeting to them. 2 Tim. 4:19.
A chief city in the Moabite territory. In Jerome's time it was called Areopolis. It is identified with Rabba,, 31 19' N, 35 38' E, about 10 miles from the Dead Sea. Num. 21:15, 28; Isa. 15:1. In other passages the name Ar appears to include the land of the Moabites. Deut. 2:9, 18, 29.
Son of Jether, of the tribe of Asher. 1 Chr. 7:38.
City in the hill country of Judah. Joshua 15:52. Identified with er-Rabiyeh, 31 26' N, 35 2' E.
This occurs as a proper name only once in the A.V. where it should read 'the Arabah,' Joshua 18:18; but it occurs in many other passages where it is translated 'a plain' or 'the plain,' and is also translated 'desert,' 'wilderness,' etc. It refers to the plain situated between two series of hills that run from the slopes of Hermon in the north to the Gulf of Akaba in the far south. It is in this plain that the Jordan runs, and in which is the Sea of Galilee and the Dead Sea, also called 'the Sea of the Plain.' About 7 miles south of the Dead Sea the plain is crossed by some hills: all north of this is now called el-Ghor, but the plain south of it retains the name of the Wady-el-Arabah. This latter part is about 100 miles in length, and the northern part about 150, so that for nearly 250 miles this wonderful plain or valley extends.
It might naturally be thought that the Jordan had at some time, after running into the Dead Sea, continued to run south until it poured itself into the Gulf of Akaba. But this is not probable, for the Dead Sea is nearly 1,300 feet below the sea, and the southern part is from end to end higher than the Ghor, The width of the Arabah is in some parts about 15 miles, but further south not more than 3 or 4. The southern end is also called the Wilderness of Zin, and it was in this part of the Arabah that a good deal of the wanderings of the people of Israel took place, before they turned to the east and left the plain on their left.
There can be no doubt that scripture uses the name 'Arabah' for the whole of the plain, both north and south. The northern part is referred to in Deut. 3:17; Deut. 4:49; Joshua 3:16; Joshua 12:3; Joshua 18:18: and the southern part in Deut. 1:1; Deut. 2:8. In other passages, especially in the prophetic books, the plain in general may be alluded to. It extends nearly due north and south, but bears toward the west before it reaches the Gulf.
A very large country is embraced by this name, lying south, south-east, and east of Palestine. It was of old, as it is now by the natives, divided into three districts.
1. Arabia Proper, being the same as the ancient Arabia Felix, embraces the peninsula which extends southward to the Arabian Sea and northward to the desert.
2. Western Arabia, the same as the ancient Arabia Petraea, embraces Sinai and the desert of Petra, extending from Egypt and the Red Sea to about Petra.
3. Northern Arabia, which joins Western Arabia and extends northward to the Euphrates.
1 Kings 10:15; 2 Chr. 9:14; Isa. 21:13; Jer. 25:24; Ezek. 27:21; Gal. 1:17; Gal. 4:25. See ARABIANS.
We read that Abraham sent the sons of Keturah and of his concubines "eastward, to the east country." Gen. 25:6. There were also the descendants of Ishmael and those of Esau. Many of these became 'princes,' and there can be no doubt that their descendants still hold the land. There are some who call themselves Ishmaelite Arabs, and in the south there are still Joktanite Arabs. We read of Solomon receiving gifts or tribute from the kings of Arabia. 1 Kings 10:15. So did Jehoshaphat, 2 Chr. 17:11 ; but in the days of Jehoram they attacked him, plundered his house, and carried away his wives and some of his sons, 2 Chr. 21:17; 2 Chr. 22:1. They were defeated by Uzziah. 2 Chr. 26:7.
During the captivity some Arabians became settlers in Palestine and were enemies to Nehemiah. Cf. Neh. 2:19; Neh. 4:7; Neh. 6:1. Among the nations that had relations with Israel, and against whom judgement is pronounced are the Arabians. Isa. 21:13-17; Jer. 25:24. And doubtless they will be included in the confederacies that will be raised against God's ancient people when Israel is again restored to their land. Cf. Ps. 83.
In the N.T. 'Arabians' were present on the day of Pentecost, but whether they were Jews or proselytes is not stated. Acts 2:11.
1. A royal city of the Canaanites, in the south, near Mount Hor, whose king fought against Israel, but who was by the help of God destroyed, both he and his people. Num. 21:1-3; Num. 33:40; Joshua 12:14; Judges 1:16. (In the two passages in Numbers read 'the Canaanite king of Arad.') It is identified with Tell Arad, 31 17' N, 35 7' E.
2. Son of Beriah, a descendant of Benjamin. 1 Chr. 8:15.
1. Son of Ulla, a descendant of Asher. 1 Chr. 7:39.
2. Father of a family who returned from exile. Ezra 2:5; Neh. 7:10.
3. A Jew whose grand-daughter married Tobiah the Ammonite, who greatly hindered the building of the city Neh. 6:18.
1. Son of Shem. Gen. 10:22, 23; 1 Chr. 1:17.
2. Son of Kemuel, Abraham's nephew. Gen. 22:21.
3. Son of Shamer, of the tribe of Asher. 1 Chr. 7:34.
4. Son of Esrom, and father of Aminadab. Matt. 1:3, 4; Luke 3:33: called RAM, Ruth 4:19; 1 Chr. 2:9, 10.
5. Place in the land of Gilead, east of the Jordan, which Jair captured. 1 Chr. 2:23.
This is the name of a large district lying north of Arabia, north-east of Palestine, east of Phoenicia, south of the Taurus range, and west of the Tigris. It is generally supposed that the name points to the district as the 'Highlands,' though it may be from Aram the son of Shem, as above. The word occurs once untranslated in Num. 23:7, as 'Aram' simply, from whence Balaam was brought, 'out of the mountains of the east;' but it is mostly translated Syria or Syrian. Thus we have -
1. ARAM-DAMMESEK, 2 Sam. 8:5, translated 'Syrians of Damascus,' embracing the highlands of Damascus including the city.
2. ARAM-MAACHAH, 1 Chr. 19:6, translated 'Syria-maachah,' a district on the east of Argob and Bashan.
3. ARAM-BETH-REHOB, 2 Sam. 10:6, translated 'Syrians of Beth-rehob: cf. Judges 18:28, a district in the north, near Dan.
4. ARAM-ZOBAH, 2 Sam. 10:6, 8, translated 'Syrians of Zoba,' a district between and Damascus, but not definitely recognised.
5. ARAM-NAHARAIM signifying 'Aram of two rivers,' Gen. 24:10; Deut. 23:4; Judges 3:8; 1 Chr. 19:6, translated 'Mesopotamia.' The two rivers are the Euphrates and the Tigris. The district would be the highlands from whence the rivers issue to the plain, and the district between the two rivers without extending to the far south.
This word occurs 2 Kings 18:26; Ezra 4:7; and Isa. 36:11, where it is translated 'the Syrian language' or 'tongue;' also in Dan. 2:4, where it is 'Syriack.' Aramaic is the language of Aram, and embraces the language of Chaldee and that of Syria. Mesopotamia, Babylonia and Syria were its proper home. The first time we meet with it in scripture is in Gen. 31:47, where Laban called the heap of witness 'Jegar-sahadutha,' which is Chaldee; whereas Jacob gave it a Hebrew name, 'Galeed.' In 2 Kings 18:26; Isa. 36:11 the heads of the people asked Rab-shakeh to speak to them in Aramaic that the uneducated might not understand what was said. In Ezra 4:7 the letter sent to Artaxerxes was written in Aramaic, and interpreted in Aramaic, that is, the copy of the letter and what follows as far as Ezra 6:18 is in that language and not in Hebrew. So also is Ezra 7:12-26.
In Daniel 2:4 the Chaldeans spoke to the king in Aramaic, the popular language of Babylon, and what follows to the end of chap. 7: is in that language, though commonly called Chaldee. This must not be confounded with the 'learning and the tongue of the Chaldeans' in Dan. 1:4, which is the Aryan dialect and literature of the Chaldeans, and probably the ordinary language which Daniel spoke in the court of Babylon. Jer. 10:11 is a verse in Aramaic.
This language differs from the Hebrew in that it avoids the sibilants. Where the Hebrew has ז z, שׁ sh, צ tz, the Aramaic has ד d, ת th, and ט t. Letters of the same organ are also interchanged, the Aramaic choosing the rough harder sounds. The latter has fewer vowels, with many variations in the conjugation of verbs, etc.
When the ten tribes were carried away, the colonists, who took their place, brought the Aramaic language with them. The Jews also who returned from Babylon brought many words of the same language. And, though it doubtless underwent various changes, this was the language commonly spoken in Palestine when our Lord was on earth, and is the language called HEBREW in the N.T., and is the same as the Chaldee of the Targums. In the ninth century the language in Palestine gave way to the Arabic, and now Aramaic is a living tongue only among the Syrian Christians in the district around Mosul.
A female belonging to Aram. 1 Chr. 7:14.
Descendant of Seir the Horite. Gen. 36:28; 1 Chr. 1:42.
A kingdom which was called upon by God, in conjunction with Medes, Persians, and others, under one captain, Cyrus, to punish Babylon in revenge of Israel. Jer. 51:27. It is identified with Urartu or Urardhu of the Assyrian inscriptions, a district in Armenia, in which is Mount Ararat, on some part of which the ark of Noah rested. Gen. 8:4. The mount is situate 39 45' N, 44 28' E, and its extreme height is about 17,000 feet above the sea, covered with perpetual snow. Objection has been taken to its great height, but it may not have been on its highest part that the ark rested.
The Jebusite from whom David purchased the place on which to build the altar of the Lord. 2 Sam. 24:16-24. Called ORNAN in 1 Chr. 21:15-28. In Samuel it is stated that David bought the threshing floor and the oxen for fifty shekels of silver. He there built an altar, and offered burnt offerings and peace offerings, without anything being said of his building a house for the Lord on the spot: whereas in Chronicles David gave to Ornan 600 shekels of gold by weight for the place. In 2 Chr. 3:1, 2 we learn that the threshing floor was on Mount Moriah, and that the site was prepared by David for the temple, which was built by Solomon. Doubtless therefore 'the place' included a much larger area than was needed for David's altar, and perhaps included the homestead of Araunah. This no doubt formed a part of what is now called the Temple area, or Mosque enclosure, in the S.E. of Jerusalem, but on what part of that area the temple was built is not known.
Arba, Arbah. [Ar'ba, Ar'bah]
Father of Anak, head of the Anakim, who were also giants. Num. 13:33. Their city was Hebron. Gen. 35:27; Joshua 14:15; Joshua 15:13; Joshua 21:11. The 'city of Arba' is elsewhere called KIRJATH-ARBA, which was afterwards called HEBRON.
Native of the northern Arabah, or el-Ghor. 2 Sam. 23:31; Chr. 11:32.
Designation of Paarai, one of David's mighty men. 2 Sam. 23:35.
The word elam occurs only in Ezek. 40:21-36, and in the A.V. is translated 'arch;' but this is judged not to be its meaning, though it is not at all certain as to what it really refers. In the margin it reads, 'galleries' or 'porches,' elsewhere 'vestibule,' and again 'projection.'
Son of Herod the Great by Malthace, a Samaritan. He succeeded his father as Ethnarch of Idumea, Judaea, Samaria, and the maritime cities of Palestine. From his known oppressive character Joseph feared to bring back the infant Jesus into his territory, and turned aside to Galilee, which was under the jurisdiction of his brother Antipas. Matt. 2:22. He reigned 10 years. Josephus relates that soon after his accession he put to death 3,000 Jews: eventually, for his tyranny to the Jews and the Samaritans he was deposed and banished to Vienne in Gaul.
People removed from Assyria to Samaria. They joined in the petition to Artaxerxes against the Jews. Ezra 4:9. The origin of the name is unknown.
City on the border of Ephraim. Joshua 16:2. Identified with Ain Arik, 31 54' N, 35 8' E.
A Christian teacher at Colosse, whom Paul calls his fellow soldier, and exhorts to fulfil his ministry. Col. 4:17; Philemon 2.
The designation of Hushai, David's friend. 2 Sam. 15:32; 2 Sam. 16:16; 2 Sam. 17:5, 14; 1 Chr. 27:33.
The word ash or aish has always been a difficult one to translate, the versions differing much; but it is now pretty well agreed that the allusion is not to the star known as Arcturus, but to the constellation known as the Great Bear; 'his sons' are supposed to be the stars in the tail of the bear. In the northern hemisphere this constellation is seen all the year round, with its apparent ceaseless motion around the north star, which none but the mighty God can guide. Job 9:9; Job 38:32. It is translated 'the Bear' in the R.V.
1. Son of Benjamin. Gen. 46:21.
2. Son of Bela, son of Benjamin (called ADDAR in 1 Chr. 8:3), whose descendants are ARDITES. Num. 26:40.
Son of Caleb, son of Hezron. 1 Chr. 2:18.
Areli, Arelites. [Are'li, Are'lites]
Son of Gad, and his descendants. Gen. 46:16; Num. 26:17.
One connected with the court of Areopagus at Athens, where Dionysius heard Paul and "clave to him and believed." Acts 17:34.
Areopagus, or Mars Hill. [Areop'agus, or Mars' Hill]
The hill of Ares, or Mars. Here was held the highest and most ancient and venerable court of justice in Athens for moral and political matters. It was composed of those who had held the office of Archon unless expelled for misconduct. Paul, who had been disputing daily in the market place, was conducted by some of the Epicurean and Stoic philosophers to Mars' Hill, not for any judicial purpose, but doubtless that they might hear him more quietly. Here he delivered his address respecting God, so suited to the heathen philosophers who heard him, and which was not without its fruit. Acts 17:19. The Greek words are Areios-pagos, but are translated Mars' Hill in Acts 17:22. The court was situate on a rocky hill opposite the west end of the Acropolis. Sixteen stone steps still lead up to the spot.
The common appellation (like Pharaoh for Egyptian kings) of the Arabian kings of the northern part of Arabia. The deputy of Aretas in Damascus sought to arrest Paul. 2 Cor. 11:32. This king, who was father-in-law to Herod Antipas, made war against him for divorcing his daughter, and defeated him. Vitellius, governor of Syria was ordered to take Aretas dead or alive; but Tiberius died before this was accomplished. Caligula, who succeeded to the empire, banished Antipas. He made certain changes in the East, and it is supposed that Damascus was detached from the province of Syria and given to Aretas.
1. A district lying to the south of Damascus and which formed a part of Bashan, where the giants resided. It had at one time 60 cities, which were ruled over by Og. Its name signifies 'stony' and it forms a remarkable plateau of basalt, which rises some 30 feet above the surrounding fertile plain, and extends 22 miles N. and S. and 14 miles E. and W., the boundary line being marked by the Bible word chebel, which signifies 'as by a rope.' Og was conquered by Moses, and Jair of Manasseh took the fortified cities, and it became a part of Manasseh's lot. Later it was called Trachonitis, and is now known as el-Lejah. There are many houses still in the district which, because of their massive proportions, are supposed to have been built by the giants. Deut. 3:3, 4, 13, 14; 1 Kings 4:13.
2. One, apparently in the service of Pekahiah, killed by Pekah. 2 Kings 15:25.
Son of Haman, slain and hanged. Esther 9:9.
Son of Haman, slain and hanged. Esther 9:8.
One, apparently in the service of Pekahiah, killed by Pekah. 2 Kings 15:25.
1. Symbolical name of Jerusalem, signifying 'Lion of God,' probably in reference to the lion being the emblem of Judah. Isa. 29:1, 2, 7. In the margin of Ezek. 43:15, the altar is called the 'lion of God;' but the word is slightly different and is translated by some the 'hearth of God,' the place for offering all sacrifices to God.
2. One whom Ezra sent to Iddo at Casiphia. Ezra 8:16.
3. In 2 Sam. 23:20; 1 Chr. 11:22, we read that Benaiah slow two 'lion-like men,' which some prefer to translate 'two [sons] of Ariel.' The Hebrew is literally 'two lions of God.'
The city of Joseph, the 'honourable counsellor,' who was permitted by Pilate to take down the body of the Lord and bury it in his own new to tomb. Matt. 27:57; Mark 15:43; Luke 23:51; John 19:38. It has not been identified, but has been supposed to be the same as Ramah, the birth-place of Samuel.
1. King of Ellasar in the East. Gen. 14:1, 9.
2. Captain of Nebuchadnezzar's guard. Dan. 2:14, 15, 24, 25.
Son of Haman the Agagite, slain and hanged. Esther 9:9.
A Macedonian of Thessalonica, companion of Paul on several journeys and on his way to Rome. Paul once calls him 'my fellow prisoner.' Acts 19:29; Acts 20:4; Acts 27:2; Col. 4:10; Philemon 24.
A resident at Rome whose household Paul saluted Rom. 16:10.
Ark of God.
This is also called 'ARK OF THE COVENANT,' 'ARK OF THE TESTIMONY,' 'ARK OF JEHOVAH.' The sacred chest belonging to the Tabernacle and the Temple. It was made of shittim wood, overlaid within and without with pure gold. It was 2-1/2 cubits long, 1-1/2 cubits in breadth, and the same in height, with a crown or cornice of gold. On each side were rings of gold in which were inserted the staves by which it was carried. Its lid, on which were the two cherubim made wholly of gold, was called the MERCY-SEAT, q.v. The ark was typical of Christ, in that it figured the manifestation of divine righteousness (gold) in man; the mercy-seat was Jehovah's throne, the place of His dwelling on earth. In the ark were placed the two tables of stone (the righteousness demanded by God from man), and afterwards the golden pot that had manna, and Aaron's rod that budded. For the place of the ark and the manner of its being moved see the TABERNACLE.
In the first journey of the children of Israel from Mount Sinai the ark of the covenant went before them to "search out a resting place for them," type of God's tender care for them. When the ark set forward Moses said, "Rise up, Lord, and let thine enemies be scattered;" and when it rested he said, "Return, O Lord, unto the many thousands of Israel." Num. 10:33-36. When they arrived at Jordan, the ark was carried by the priests 2000 cubits in front of the host that they might know the way they must go, Joshua 3:3, 4, and the ark remained on the shoulders of the priests in the bed of the river, until all had passed over. Joshua 3:17. This typifies association with Christ's death and resurrection.
The ark accompanied them in their first victory: it was carried by the priests around Jericho. It is only in the power of Christ in resurrection that the saint can be victorious. The tabernacle was set up at Shiloh, and doubtless the ark was placed therein, Joshua 18:1, though it may have been carried elsewhere. In Eli's days when Israel was defeated they fetched the ark from Shiloh that it might save them, but they were again defeated, and the ark, in which they had placed their confidence instead of in Jehovah, was seized by the Philistines. 1 Sam. 5:1. When put into the house of their god Dagon the idol fell down before it on two occasions, and on the second was broken to pieces. Subsequently it was taken from Ashdod to Gath, and from Gath to Ekron, and the people were smitten by the hand of God in each city.
After seven months a new cart was made, to which two milch kine were yoked, and the ark sent back to the Israelites with a trespass offering to the God of Israel. The kine, contrary to nature, went away from their calves, and went direct to Beth-shemesh, for it was God who restored the ark. There God smote the men of the place for looking into the ark. It was then taken to Kirjath-jearim and placed in the house of Abinadab. 1 Sam. 6; 1 Sam. 7:1, 2. See ABINADAB.
In after years David fetched the ark from thence on a new cart, but the ark being shaken, Uzzah put forth his hand to steady it, and was smitten of God. This frightened David and the ark was carried aside to the house of Obed-edom. The law had directed how the ark was to be carried, and the new cart was following the example of the Philistines: Uzzah disregarded God's plain direction and heeded not the sacredness of that which represented the presence of God. David however, hearing that God had blessed the house of Obed-edom, again went for the ark, and now it was carried by the Levites according to divine order, and with sacrifices and rejoicing it was placed in the tabernacle or tent that David had pitched for it. 2 Sam. 6.
When Solomon had built the temple, the ark was removed thither, and the staves by which it had been carried were taken out: the ark had now found its resting place in the kingdom of Solomon, whose reign is typical of the millennium. It is significant too that now there were only the two tables of stone in the ark, 1 Kings 8:1-11: the manna had ceased when they ate of the old corn of the land, which is typical of a heavenly Christ; and the witness of Aaron's rod was no longer needed now they were in the kingdom. The wilderness circumstances, in which the manna and the priesthood of Christ were so necessary, were now passed. These are both mentioned in Heb. 9:4, for there the tabernacle, and not the temple is in contemplation.
No further mention is made of the ark: it is supposed to have been carried away with the sacred vessels to Babylon, and to have never been returned: if so there was no ark in the second temple nor in the temple built by Herod, nor do we read of the ark in connection with the temple described by Ezekiel. In Rev. 11:19 the ark of God's covenant is seen in the temple of God in heaven: symbol here of the resumption of God's dealings with His earthly people Israel.
Ark of Noah.
The vessel constructed by the command of God, by which Noah and his household and some of every living creature of the earth were saved when the world was destroyed by the flood. Precise instructions were given by God as to the construction of the ark. It was to be made of 'gopher' wood, a kind known at the time, but which cannot now be identified with certainty; and it was to be pitched within and without with pitch, or bitumen, to make it water-tight.
Its proportions were to be 300 cubits long, 50 cubits broad, and 30 cubits high. If the cubit be taken at 18 inches, its length would have been 450 feet, its breadth 75 feet and its height 45 feet. If the cubit used had been 21 inches, the dimensions would be one-sixth larger.
A window was to be made to the ark. Gen. 6:16. The word tsohar signifies 'a place of light' and was probably placed in the roof, and may have served in some way for ventilation as well as for giving light. Another word for window is used in Gen. 8:6 (challon) which could be opened from the inside. This word is used for the windows or casements of houses, and would give ventilation. In Gen. 6:16, after speaking of the window, it says, "and in a cubit shalt thou finish it above;" it is a question whether this refers to the size of the window or whether the word 'it' refers to the ark. It has been said that the feminine suffix, which is rendered 'it' cannot refer to the word window, which is masculine: so that it is possible the cubit refers to the roof; that the middle of the roof should be raised, giving a cubit for the pitch of the roof. A door was to be made in the side of the ark; and the ark was to be divided into three stories. 'Rooms,' or 'nests' (margin) are also mentioned. Gen. 6:14.
Such is the description given us of the form of the ark. It was by faith Noah prepared the ark, by which he condemned the world, and became heir of the righteousness which is by faith. Heb. 11:7. It is thus referred to in 1 Peter 3:20, 21, "into which few, that is, eight souls, were saved through water: which figure also now saves you, [even] baptism, not a putting away of [the] filth of flesh, but [the] demand as before God of a good conscience, by [the] resurrection of Jesus Christ."
It may just be added that the form of the ark was not intended for navigation amid storms and billows, but it was exactly suited for the purpose for which it was constructed. A ship for freight was once made in like proportions, to be used in quiet waters, and was declared to be a great success.
Various questions have been raised as to the veracity of the Bible account of the Deluge, for which see FLOOD.
Ark of Bulrushes.
The little boat or cradle in which Moses was placed by his mother. It was made of bulrushes, or rather paper-reeds or papyrus which grew in the river Nile. It was daubed with slime and with pitch, that is, most probably first covered with wet earth or clay, and then with bitumen. Ex. 2:3, 5. Some of the heathen writers speak of the papyrus-woven craft of the Nile. God answered the faith of the parents, and Moses was drawn out of the water to be the saviour of His people.
Tribe descended from Canaan, son of Ham; it probably resided in Arca, in the north of Phoenicia, about 15 miles north of Tripoli, now called Tell Arka. Gen. 10:17; 1 Chr. 1:15.
The member of the body which is capable of lifting burdens and defending the person: it is used symbolically for the power and strength of God on behalf of His saints. Ex. 15:16; Ps. 77:15; Isa. 51:9; Isa. 53:1. The arm of Jehovah is often spoken of in the O.T. It redeemed, Ex. 6:6; etc.; gathers His own, Isa. 40:11; and rules for Him, Isa. 40:10, as in the kingdom. It is a holy arm, Isa. 52:10; Ps. 98:1; and it is a glorious arm, Isa. 63:12. The arm of the Lord is revealed to souls where there is repentance and faith in the report which God sends. Isa. 53:1; Rom. 10:16. It is to be trusted in even by the isles of the Gentiles, that is, by sinners everywhere in creation. Isa. 51:5.
The Hebrew name of the place where the kings of the earth and of the whole world will be gathered together to make war against the Lord Jesus in the great day of Almighty God. Rev. 16:16. There seems to be an allusion to the great battle field of Palestine in the Esdraelon, and to the Megiddo mentioned in Judges 5:19; 1 Kings 4:12; 2 Kings 23:29, 30. The word itself is translated 'the mountain of slaughter,' and may be used symbolically for the destruction that will surely fall upon the enemies of the Lord Jesus.
This name occurs in the A.V. in 2 Kings 19:37; Isa. 37:38, as the place to which two sons of Sennacherib fled after killing their father; but in both these passages the Hebrew word is Ararat. Armenia occurs in the LXX in the passage in Isaiah. Armenia lies west of the Caspian Sea, and extends northward of 38 N. lat. It is now partly in the Russian and partly in the Turkish empires.
Son of Saul and Rizpah, hanged by the Gibeonites. 2 Sam. 21:8.
None of the Hebrew words translated 'armour' refer definitely to what is understood now by armour worn on the person. Saul armed David with his 'armour,' 1 Sam. 17:38, but the word used is also translated 'clothes,' etc., and it may refer to Saul's warrior-dress. The articles named are somewhat more definite.
1. Saul put on David a 'HELMET of brass.' These were raised a little above the head, as may be seen by some of the sculptures from Nineveh. 1 Sam. 17:38; Ezek. 23:24: the word is qoba. Another word, koba, meaning the same, is found in 1 Sam. 17:5; 2 Chr. 26:14; Isa. 59:17; Jer. 46:4; Ezek. 27:10; Ezek. 38:5.
2. COAT OF MAIL. Saul put on David a 'Coat of Mail,' shiryon. 1 Sam. 17:5, 38. This word is translated 'HABERGEON ' in 2 Chr. 26:14 ; Neh 4:16, which also signifies 'coat of mail,' and there is a similar word in Job 41:26. It was made of brass scales fastened together. The weight of Goliath's coat of mail was 5,000 shekels.
3. GREAVES. The giant wore Greaves of brass upon his legs. 1 Sam. 17:6. The word is mitschah, and occurs nowhere else.
4. TARGET. He had a Target of brass between his shoulders, 1 Sam. 17:6: the word is kidon, and is elsewhere translated both 'shield' and 'spear.' In this case it was probably a small spear carried between the shoulders.
5. SHIELD. A Shield was carried before him. This was a tsinnah, a shield of large size to protect the whole body, with a large boss in the centre rising to a point which could be used as a weapon. It is employed figuratively for God's protecting care of His people. Ps. 5:12; Ps. 91:4. The same word is translated BUCKLER. Ps. 35:2; Ezek. 23:24; Ezek. 26:8, etc.
Another word is used for a smaller shield, magen, and this is the word which occurs most commonly in the O.T., especially in the Psalms, referring to God's protection, as Ps. 28:7; Ps. 33:20; Ps. 84:11; Ps. 119:114, etc. The same word is translated BUCKLER. 2 Sam. 22:31; 1 Chr. 5:18; Cant. 4:4; Jer. 46:3, etc.
The word shelet is translated Shield, but is also applied to Shields of gold, 2 Sam. 8:7, and those suspended for ornament. Ezek. 27:11. It occurs also in 2 Kings 11:10; 1 Chr. 18:7; 2 Chr. 23:9; Cant. 4:4; Jer. 51:11.
In the N.T. 'armour' is used symbolically.
1. ὅπλα in contrast to 'the works of darkness' we are exhorted to put on 'the armour of light.' Rom. 13:12. Paul and his fellow-labourers commended themselves as God's ministers by the "armour, or arms, of righteousness on the right hand and on the left." 2 Cor. 6:7.
2. πανοπλία, 'whole armour.' One stronger than Satan takes away all his 'armour.' Luke 11:22. The Christian is exhorted to put on the 'whole armour of God,' the panoply, that he may stand in the evil day in his conflict with the spiritual powers of wickedness in the heavenlies. Eph. 6:11, 13. See BREASTPLATE, HELMET, etc.
An attendant on a warrior, filling a place of trust and honour. When Saul loved David he made him his armourbearer. 1 Sam. 16:21. On Saul being wounded, his armourbearer refused to kill him; but when Saul was dead the armourbearer fell upon his sword and died also. 1 Sam. 31:5.
In Neh. 3:19 the word is nesheq also translated 'armour.' In Cant. 4:4 it is talpiyyoth, 'armoury' or heap of swords. In Jer. 50:25 it is otsar, signifying 'treasury.'
The offensive arms found in the O.T. are:
1. The SWORD, for which several Hebrew words are used: a. baraq, often translated 'lightning;' it is 'glittering sword' in Job 20:25. b. chereb, a sword, as laying waste. It is the word commonly used in the O.T. for sword (everywhere indeed except in the references given here under the other words): it was a straight tapering weapon, with two edges and a sharp point. Ps. 149:6; Isa. 14:19. It is used metaphorically for keen and piercing words, as in Ps. 57:4; Ps. 64:3. c. retsach, an undefined slaying weapon, translated 'sword' only in Ps. 42:10. d. shelach, a missile of death, as a dart. Job 33:18; Job 36:12; Joel 2:8. e. pethichoth, from 'to open,' is translated 'drawn sword' in Ps. 55:21.
2. SPEARS. a. chanith, thus named as being flexible: it is the word mostly used for the spear. 1 Sam. 13:19; Ps. 57:4. It is this weapon that will be beaten into pruning hooks. Isa. 2:4; Micah 4:3. b. kidon, a smaller kind of lance, or javelin. Joshua 8:18, 26; Job 41:29; Jer. 6:23. c. tselatsal, harpoon. Job 41:7. d. qayin, lance, 2 Sam. 21:16. e. romach, spear used by heavy-armed troops, the iron head of a spear. Judges 5:8, etc. The pruning hooks are to be beaten into spears in the time of God's judgements. Joel 3:10.
3. BOW, from which arrows are discharged, qesheth, generally made of wood, but sometimes of steel or brass. Job 20:24. It is constantly found in the O.T. from Genesis to Zechariah. It is used to express punishment from God, Lam. 2:4; Lam. 3:12; and of men to show their power to injure. Ps. 37:14, 15. 'A deceitful bow' expresses a man who fails just when his aid is most needed, as when a bow breaks suddenly. Ps. 78:57; Hosea 7:16.
4. The SLING, by which stones are discharged, qela. It was by means of this that David smote Goliath. 1 Sam. 17:40, 49, 50. Of the Benjamites there were 700 men lefthanded; "every one could sling stones at an hair breadth, and not miss." Judges 20:16. (In Prov. 26:8 occurs another word for sling margemah, but the passage is considered better translated "as he that putteth a precious stone in a heap of stones," as in the margin.)
5. 'ENGINES,' with which Uzziah shot arrows and great stones. 2 Chr. 26:15.
It must be remembered that Israel were the hosts of Jehovah, keeping His charge and fighting His battles. Ex. 12:41; Joshua 5:14. It appears that all who reached the age of twenty years were contemplated as able to bear arms, Num. 1:3; and they marched and encamped in 4 divisions of 3 tribes each, with a captain over every tribe. The subdivisions were into thousands and hundreds, Num. 31:14, and into families. Joshua 7:17. There were also trumpet calls, Num. 10:9 (cf. 1 Cor. 14:8), and all the appearance of careful organisation. Until the time of the kings this natural or tribal organisation seems to have been usual, but in the time of Saul there was a body guard, 1 Sam. 13:2, and a captain of the host, 1 Sam. 17:55. In David's days those heroes who were with him in the cave of Adullam formed the nucleus of his 'mighty men.' 2 Sam. 23:8-39. They were devoted to the service of God's king. David afterwards organised a monthly militia of 24,000 man under 12 captains. 1 Chr. 27:1-15.
The general gradation of ranks was into privates; 'men of war;' officers; Solomon's 'servants;' captains or 'princes;' and others variously described as head captains, or knights or staff officers; with rulers of his chariots and his horsemen. 1 Kings 9:22. It may be noticed that horses having been forbidden, Deut. 17:16, it was not until Solomon's time that this was organised, though David had reserved horses for a hundred chariots from the spoil of the Syrians. 2 Sam. 8:4. Solomon, trading with Egypt, 1 Kings 10:28, 29, enlarged their number until the force amounted to 1,400 chariots, and 12,000 horsemen, 1 Kings 10:26; 2 Chr. 1:14. Every able man being a soldier gave David the immense army of 1,570,000 men that 'drew sword.' 1 Chr. 21:5. After the division, Judah under Abijah had an army of 400,000 'valiant men,' and Israel at the same time of 800,000 'chosen men.' Afterwards Asa had 580,000 'mighty men of valour;' and Jehoshaphat, who had waxed great exceedingly, had as many as 1,160,000 men, besides those left in the fenced cities. 2 Chr. 17:14-19.
In the N.T. a few references are made to the Roman army. A 'Legion' was a body that contained within itself all the gradations of the army. It might be called under the empire, in round numbers, a force of not more than 6,000 men. Every legion at times contained 10 cohorts of 600 each; every cohort 3 maniples of 200; and every maniple 2 centuries of 100: hence the name of centurion or commander of 100 men, as found in Acts 10:1, 22, etc. Each legion was presided over by 6 chiefs, χιλίαρθος, each commanding 1,000 men, mostly translated 'chief captain,' as in Acts 21:31-37, etc.: it is 'high captain' in Mark 6:21; and 'captain' in John 18:12; Rev. 19:18. A cohort, σπεῖρα, is translated 'band' in Acts 10:1; Acts 21:31, etc. A 'quaternion' embraced 4 soldiers. Acts 12:4.
The head quarters of the Roman troops was at Caesarea, with a cohort at Jerusalem; but at the time of the feast, when, alas, the mutinous disposition of the Jews was sure to appear, additional troops were present in the city but without their standards of the eagle, etc., which were especially obnoxious to the Jews. Though the Romans were God's rod to punish them, their stiff necks could not bow, nor receive the punishment as from Jehovah.
Descendant of David. 1 Chr. 3:21.
Ravine or wady with its mountain torrent, which formed the border between Moab and Ammon, now known as Wady Mojib. It has sources both north and south which unite, and its stream running nearly east and west, rushes through a deep ravine and falls into the Dead Sea at about its centre north and south. Num. 21:13-28; Num. 22:36; Deut. 2:24, 36; Judges 11:13-26; Isa. 16:2; Jer. 48:20; etc.
Arod, Arodi, Arodites. [A'rod, Aro'di, Aro'dites]
Son of Gad, and his descendants. Gen. 46:16; Num. 26:17.
1. City 'before Rabbah,' that is, near Rabbath Ammon, in the valley of the Jabbok, built or rebuilt by the tribe of Gad. Num. 32:34; Joshua 13:25; 2 Sam. 24:5.
2. Moabite city on the north bank of the Arnon. Deut. 2:36; Joshua 13:9, 16; Judges 11:26; 2 Kings 10:33. Identified with Arair, 31 27' N, 35 43' E.
3. District near Damascus. Isa. 17:2.
4. City in Judah, S.E. of Beersheba. 1 Sam. 30:28. Identified with Ararah, 31 11' N, 34 56' E.
Designation of Hothan, father of two of David's captains. 1 Chr. 11:44.
Arpad, Arphad. [Ar'pad, Ar'phad]
Fortified city near Hamath. 2 Kings 18:34; 2 Kings 19:13; Isa. 10:9; Isa. 36:19; Isa. 37:13; Jer. 49:23.
Son of Shem, born two years after the flood, from whom Abraham descended. Gen. 10:22, 24; Gen. 11:10-13; 1 Chr. 1:17,18, 24. Stated as the father of Cainan in Luke 3:36. See CAINAN.
With the bow, a common weapon of the ancients. We know not of what wood the arrows of the Israelites were made. Apparently the arrows were sometimes poisoned. Job 6:4; Ps. 120:4; Num. 24:8; Deut. 32:23, etc. Arrows are used metaphorically for the judgements of God, Ps. 38:2; Ps. 45:5: also for anything sharp and painful, as smiting by the tongue. Jer. 9:8.
1. Persian king, identified as the magian impostor who pretended to be Smerdis the brother of Cambyses. When appealed to by the adversaries of the Jews, he stopped the building of the temple. He was slain after a reign of eight months. Ezra 4:7, 8, 11, 23.
2. Another Persian king identified as Artaxerxes Longimanus B.C. 474-434, son of Xerxes, the Ahasuerus of Esther. He greatly favoured both Ezra and Nehemiah; he beautified the temple or bore the expense of its being done, Ezra 7:27, and under his protection the wall of the city was finished. Ezra 6:14; Ezra 7:1-21; Ezra 8:1; Neh. 2:1; Neh. 5:14; Neh. 13:6. It was in the 20th year of this king that the command to build the city was given, from which began the dates of the prophecy of the Seventy weeks of Daniel, which is fixed by Usher and Hengstenburg at B.C. 454-5. For the succession of the Persian kings see PERSIA.
Companion of Paul at Nicopolis. Titus 3:12.
Name of the heathen goddess Diana, as given in the Greek of Acts 19:24-35: she was regarded as presiding over the productive and nutritive powers of nature.
A general name for skilled artisans, whether in metal, stone, or wood. Tubal-cain was the first named as an artificer in brass and iron. Jubal was the father of all such as handled, or invented and made, the harp and the organ. Cain also built a city. Gen. 4:17, 21, 22. In the above we see the application of the arts by man at a distance from God to promote their own welfare in independence of God. In after times the spirit of wisdom was given to Bezaleel for the work of the tabernacle in "all manner of workmanship." Ex. 35:31: cf. also 1 Chr. 29:5; 2 Chr. 34:11. It would seem that the Jews never afterwards lost this skill, as the remains of the walls of Jerusalem indicate. Nebuchadnezzar carried off all the craftsmen (same word as artificers) and smiths from Jerusalem, 2 Kings 24:14, and he may have made use of their skill to adorn Babylon.
A general term for tools, armour, etc. In 1 Sam. 20:40 it refers to the bow and arrows Jonathan had used.
The third commissariat district of Solomon, probably the rich corn-growing country in the Shephelah or low hills of Judah. 1 Kings 4:10.
City or district apparently near Shechem, the abode of Abimelech. Judges 9:41. Identified with el-Ormeh, 32 9' N, 35 19' E.
Island on the Phoenician coast: now called Ruad, about 34 51' N, 35 52' E. Ezek. 27:8, 11.
Family name of one of the sons of Canaan. Gen. 10:18; 1 Chr. 1:16: doubtless connected with the island of Arvad.
Steward of Elah, king of Israel. 1 Kings 16:9.
1. Great grandson of Solomon and king of Judah, B.C. 955-914. "Asa did that which was right in the sight of the Lord, as did David his father." He removed the idols his fathers had made, 1 Kings 15:11, and he deposed Maachah, his mother, or perhaps grandmother, from being queen because she favoured idolatry. On the country being invaded by the Ethiopians with a million troops and 300 chariots, he cried to the Lord, who fought for him, and the enemy was smitten. He was counselled by Azariah not to forsake the Lord, which led to the spoil being offered to God, and to the king and his people entering into a covenant to seek the Lord.
Subsequently Asa was threatened by Baasha king of Israel who began to build Ramah, a fortified city only a few miles from Jerusalem. To stop this Asa paid a large sum of money to Benhadad king of Syria to invade Israel. This was for the time successful: the building of Ramah was stopped, and Asa carried away the stones thereof and built Geba and Mizpah.
This recourse for aid to the king of Syria, who was an idolater, was very displeasing to God, and the king was rebuked by Hanani the seer. While Asa trusted in the Lord he had deliverance, but having relied on the king of Syria, he should have war all his days. Asa, alas, did not humble himself, but put Hanani in prison, and oppressed some of the people. He was disciplined in his person, for he was diseased in his feet, and the disease increased exceedingly; yet he sought not the Lord, but to the physicians (perhaps these were healers by magic arts in connection with idolatry, on which God's blessing could not be asked) and he died after a reign of 41 years. 1 Kings 15.; 2 Chr. 14, 15, 16.; Matt. 1:7, 8.
2. A Levite, the father of Berechiah. 1 Chr. 9:16.
1. Nephew of David, being son of his sister Zeruiah; he was a valiant man and one of David's captains; was slain by Abner while pursuing him. 2 Sam. 2:18-32; 2 Sam. 3:27, 30; 1 Chr. 11:26; 1 Chr. 27:7.
2. Levite sent by Jehoshaphat to teach the law in the cities of Judah. 2 Chr. 17:8.
3. Levite in Hezekiah's time, an overseer of tithes, etc. 2 Chr. 31:13.
4. Father of Jonathan who returned from exile. Ezra 10:15.
Asahiah, Asaiah. [Asahi'ah, Asai'ah]
1. An officer sent by Josiah to Huldah the prophetess after the book of the law had been found. 2 Kings 22:12, 14; 2 Chr. 34:20.
2. Descendant of Simeon. 1 Chr. 4:36.
3. Descendant of Merari. 1 Chr. 6:30.
4. A Shilonite who became a dweller in Jerusalem. 1 Chr. 9:5.
5. Descendant of Merari who assisted in bringing up the ark from Obed-edom's house, 1 Chr. 15:6, 11 (possibly the same as No. 3).
1. A leader of the choir in David's time, and once called a 'seer.' 2 Chr. 29:30. He was descended from Gershom the Levite. 1 Chr. 6:39; 1 Chr. 15:17, 19; 1 Chr. 16:5, 7, 37, etc. Twelve psalms are attributed to him, namely, 50, 73 to 83. His office seems to have been hereditary. Ezra 2:41; Ezra 3:10; Neh. 7:44, etc.
2. Father of Joah recorder to Hezekiah. 2 Kings 18:18, 37; Isa. 36:3, 22.
3. A Levite, whose descendants dwelt in Jerusalem after the exile. 1 Chr. 9:15.
4. A Korhite, whose posterity were porters in the tabernacle in the time of David. 1 Chr. 26:1.
5. An officer, probably a Jew, controller of the forests of king Artaxerxes in Judaea. Neh. 2:8.
Son of Jehaleleel, a descendant of Judah. 1 Chr. 4:16.
Son of Asaph appointed by David to the service of song. 1 Chr. 25:2. Supposed by some to be the same as JESHARELAH in 1 Chr. 25:14, as noted in the margin; and by others to be the same as AZAREEL in 1 Chr. 25:18.
This term is constantly applied to the return of the Lord Jesus Christ to heaven from whence He came. John 3:13. Leading His eleven apostles out as far as Bethany, on the eastern slope of the Mount of Olives, in the act of blessing them He ascended up to heaven, and a cloud hid Him from their sight. Mark 16:19; Luke 24:50, 51; Acts 1:9. The ascension of the Lord Jesus is a momentous fact for His saints: the One who bore their sins on the cross has been received up in glory, and sits on the right hand of God.
As forerunner He has entered into heaven for the saints, and has been made a high priest for ever after the order of Melchisedec. Heb. 6:20. His ascension assured, according to His promise, the descent of the Holy Spirit, which was accomplished at Pentecost. John 16:7; Acts 1:4, 8; Acts 2:1-47. As ascended He became Head of His body the church, Eph. 1:22, and gave gifts to men, among which gifts are evangelists who preach to the world, and pastors and teachers to care for and instruct the saints. Ps. 68:18; Eph. 4:8-13.
His ascension is a demonstration through the presence of the Holy Spirit that sin is in the world and righteousness in heaven, for the very One they rejected has been received by the Father into heaven. John 16:10. The ascension is also a tremendous fact for Satan: the prince of this world has been judged who led the world to put the Lord to death; and in His ascension He led captivity captive, having broken the power of death in which men were held, Eph. 4:8, for He had in the cross spoiled principalities and powers and made a show of them openly, triumphing over them in it. Col. 2:15.
Above all, the ascension is a glorious fact for the blessed Lord Himself. Jehovah said unto Him, "Sit thou at my right hand, until I make thine enemies thy footstool." Ps. 110:1. He has taken His place as man where man never was before, and He is also glorified with the glory which He had before the world was, besides the glory which He graciously shares with His saints. John 17:5, 22.
Daughter of Poti-pherah, priest of On, wife of Joseph, and mother of Manasseh and Ephraim. Gen. 41:45, 50; Gen. 46:20.
The particular tree pointed out by the Hebrew word oren is not known. Isa. 44:14. The LXX and the Vulgate call it 'pine.'
1. Levitical city in Judah. Joshua 15:42; 1 Chr. 6:59: not identified.
2. City in Simeon. Joshua 19:7; 1 Chr. 4:32. See AIN.
A family apparently descended from Shelah who 'wrought fine linen.' 1 Chr. 4:21.
Ashbel, Ashbelites. [Ash'bel, Ash'belites]
Son of Benjamin and family descended from him. Gen. 46:21; Num. 26:38; 1 Chr. 8:1.
One of the five chief cities of the Philistines. It was assigned to Judah, but was not subdued by them, and thus became a thorn in their sides. Num. 33:55. It was to this city that the ark was taken by the Philistines, and where Dagon their fish-god fell before it. 1 Sam. 5:1-7. Uzziah broke down its wall, and built cities near it. 2 Chr. 26:6. It was on the high road from Palestine to Egypt which doubtless led Sargon king of Assyria to take it by his general, about B.C. 714. Isa. 20:1. Herodotus records that Psammetichus, king of Egypt, besieged it for 29 years. Jeremiah speaks of Ashdod as one of the places which was made to drink of the fury of God. Jer. 25:15-20. The Maccabees destroyed the city, but Gabinius rebuilt it at the time of the conquest of Judaea by the Romans, B.C. 55, and it was afterwards assigned on the death of Herod the Great to his sister Salome. It was situated about 3 miles from the Mediterranean, and midway between Gaza and Joppa. It is now called Esdud, or Esdood, 31 46' N, 34 40' E, and is wretched in the extreme, though lying in a fertile plain. It is called in the N.T. AZOTUS, where Philip was found after baptising the eunuch. Acts 8:40. Its inhabitants are referred to as ASHDODITES, ASHDOTHITES. Joshua 13:3; Neh. 4:7.
This is once translated 'springs of Pisgah,' pointing it out as a place from whence water issued, being the sides of the mountain called Pisgah, or it may apply to the range of mountains on the east of the Dead Sea, of which Pisgah was a part. Deut. 3:17; Deut. 4:49; Joshua 12:3; Joshua 13:20. It lies due east of the north end of the Dead Sea, and is now called Ayun Musa.
Asher, Aser. [Ash'er, A'ser]
Eighth son of Jacob by Zilpah, Leah's handmaid. Gen. 13. The signification of the name as in the margin is 'happy.' His posterity formed one of the twelve tribes. Its portion in the land was in the extreme north, extending northward from Mount Carmel. It was bounded on the east by Naphtali, and on the south east by Zebulon. It was doubtless intended that their west border should have been the Great Sea, but we read that they did not drive out the inhabitants of Accho, Zidon, Ahlab, Achzib, Helbah, Aphik and Rehob; but the Asherites dwelt among the Canaanites. Judges 1:31, 32. This left a tract of land on the sea coast unoccupied by Asher.
When Jacob called his sons about him to tell them what should befall them in the last days, he said of Asher, "Out of Asher his bread shall be fat, and he shall yield royal dainties." Gen. 49:20. When Moses ordained that certain of the tribes should stand on Mount Gerizim to bless the people, and certain others on Mount Ebal to curse, Asher was one of those chosen to stand on the latter. Deut. 27:13. And when Moses blessed the tribes before he died, he said of Asher, "Let Asher be blessed with children; let him be acceptable to his brethren, and let him dip his foot in oil. Thy shoes shall be iron and brass; and as thy days, so shall thy strength be." Deut. 33:24, 25.
In Jacob's prophecy as to this tribe there is depicted the future blessing of all Israel after the salvation of the Lord has come in, announced at the close of Dan's apostasy. In Deuteronomy, what is future also as to Israel, is probably presented, but connected rather with the government of God in His hands who is King in Jeshurun.
When Deborah and Barak went to the war they had to lament in their song that Asher abode by the sea coast, and came not to their aid, Judges 5:17; but when subsequently the Midianites and the Amalekites invaded the land Asher responded to the call of Gideon. Judges 6:35; Judges 7:23. At the secession of the ten tribes Asher became a part of Israel, and very little more is heard of this tribe. When Hezekiah proclaimed a solemn passover and sent invitations to the cities of Israel as well as to Judah, though many laughed the messengers to scorn, divers of Asher humbled themselves and came to Jerusalem. 2 Chr. 30:11.
When numbered at Sinai there were 41,500 able to go forth to war, and when near the promised land they were 53,400; but when the rulers of the tribes are mentioned in the time of David, Asher is omitted. Num. 1:41; Num. 26:47; 1 Chr. 27:16-22. The tribe is twice referred to in the N.T. as ASER. In Rev. 7:6, twelve thousand of Asher will be sealed, and in Luke 2:36, Anna a prophetess, of the tribe of Asher, gave thanks in the temple at the birth of the Saviour. Asher is one of the tribes still to come into blessing, and have a portion in the land. Ezek. 48:2, 3. See THE TWELVE TRIBES
One of the tribe of Asher. Judges 1:32.
Ashes, mostly from burnt wood, were used as a sign of sorrow or mourning, either put on the head, 2 Sam. 13:19, or on the body with sackcloth, Esther 4:1; Jer. 6:26; Lam. 3:16; Matt. 11:21; Luke 10:13; or strewn on a couch on which to lie, Esther 4:3; Isa. 58:5; Jonah 3:6. To eat ashes expresses great sorrow, Ps. 102:9; and to be reduced to them is a figure of complete destruction, Ezek. 28:18; Malachi 4:3; to feed on them tells of the vanities with which the soul may be occupied. Isa. 44:20. 'Dust and ashes' was the figure Abraham used of himself before Jehovah, Gen. 18:27; and Job said he had become like them by the hand of God. Job 30:19. For the ashes of the Red Heifer see HEIFER.
An idol introduced into Samaria by the colonists sent from Hamath by the king of Assyria. 2 Kings 17:30.
Ashkelon, Askelon. [Ash'kelon, As'kelon]
One of the five principal cities of the Philistines. It fell to the lot of Judah, who took Askelon and the coasts thereof, Judges 1:18, but they did not really subdue it, for it was in the hands of the Philistines when Samson, with the Spirit of the Lord upon him, slew thirty men in the city and took their spoil, Judges 14:19, and that it remained so we see from 1 Sam. 6:17, and 2 Sam. 1:20. The judgements of God were denounced against this city, Jer. 25:20; Jer. 47:5, 7; Amos 1:8; Zech. 9:5; and the remnant of Judah should dwell there. Zeph. 2:4, 7.
The city was situated on the sea coast, midway between Gaza and Ashdod: it is now called Askulan or Askalan, 31 40' N, 34 33' E. In modern times the city was held by the Crusaders, and within its walls Richard of England held his court: the walls which this king aided with his own hands to repair may, it is thought, still be traced, and masses of masonry and broken columns of granite still lie about. By the Mahometan geographers it was called the Bride of Syria.
Ashkenaz, Ashchenaz. [Ash'kenaz, Ash'chenaz]
Son of Gomer, the son of Japheth, and his descendants, who settled in the vicinity of Armenia. Gen. 10:3; 1 Chr. 1:6; Jer. 51:27.
1. Town in the west of Judah near Dan. Joshua 15:33. Identified with Hasan, 31 47' N, 34 59' E.
2. Town in the low hills of Judah, probably to the S.W. of Jerusalem. Joshua 15:43.
Prince of the eunuchs under Nebuchadnezzar. Dan. 1:3.
Descendant of Manasseh. 1 Chr. 7:14. See ASRIEL. |
Judicature Genes and Justice
The Growing Impact of the New Genetics on the Courts
November-December 1999 Vol 83(3)
COMPLEX SCIENTIFIC EVIDENCE and the JURY
by Robert D. Myers, Ronald S. Reinstein, and Gordon M. Griller
DNA—deoxyribonucleic acid, the chemical molecule inside cells which carries biological information. DNA is a double stranded molecule held together by weak hydrogen bonds between complementary base pairs of nucleotides (Adenine and Thymine; Guanine and Cytosine). This molecule carries genetic information from parent to offspring.
Genome—one copy of all the DNA found in each cell of an organism. The human genome is composed of three billion base pairs of DNA packaged as 23 chromosomes. There are two copies of each [chromosome] in a cell, one copy from each of your parents. The genome contains the organism's genes, the instructions for building that life form.
These definitions of DNA and genome, two scientific concepts at the heart of this issue of Judicature, seem rather straightforward and simple. One may think that even without scientific background and learning, these concepts can be readily understood, perhaps with a few additional definitions, or a little more explanation from someone knowledgeable. But as the twentieth century draws to a close, the U.S. Human Genome Project moves closer to its goal: determining and mapping the complete sequence of DNA in the human genome by the year 2003. The implications of the Project's work for courts and the entire legal system are enormous:
The HGP's ultimate goal is to discover all of the more than 80,000 human genes and render them accessible for further biological study.... Information obtained as part of the HGP will dramatically change almost all biological and medical research and dwarf the catalog of current genetic knowledge. Both the methods and data developed through the project are likely to benefit investigations of many other genomes, including a large number of commercially important plants and animals. In a related project to sequence the genomes of environmentally and industrially interesting microbes, in 1994 DOE initiated the Microbial Genome Program. For this reason, in addition to the DOE and NIH programs, genome research is being carried out at agencies such as the U.S. Department of AgricultureÉand the private sector. In a departure from most scientific programs, research also is being funded on the ethical, legal, and social implications (ELSI) of HGP data.1
Potential government and private sector applications of this knowledge—gene therapies, gene transfers, genetic screening, and new biotechnologies—ultimately will give rise to a myriad of disputes that will make their way into the courts for resolution. The legal issues involved in these controversies, and the evidence that underlies them, will be far more complex than the two brief definitions of DNA and genome at the outset of this article. As judges and lawyers ready themselves for this growing level of scientific evidence, one principal justice system decision maker is largely unprepared...the trial juror.
Already, the most familiar form of genomic evidence, DNA "fingerprinting" (or "profiling," or "typing") in criminal cases, is widely admissible in state and federal courts, by court decision or legislation. The possible uses of genomic evidence, however, are not limited to criminal matters. Some states have already enacted legislation regulating health insurers' use of genetic testing data. Disputes involving insurance coverage, medical malpractice, product liability, toxic torts, employment discrimination, paternity, privacy, and intellectual property will become increasingly complex as the knowledge of not only human, but plant and animal genetics, and the practical applications of that knowledge, become more widespread. As one commentator has said, it is "not whether genetic evidence will ever be admitted into court, but when and under what kinds of circumstances."2
Against this backdrop, the ability of juries to adequately understand genomic evidence, distinguish between and resolve contradicting opinions of expert witnesses, and properly apply the law to the evidence is being called into question. Some court watchers believe juries are not competent to resolve scientific evidence issues, and matters of complex scientific evidence should be removed from them. Others argue that the societal values represented by both criminal and civil juries are too important to forego, and that the common sense approach jurors bring to disputes equip them in a unique, capable manner to comprehend novel and complex scientific evidence. In reality, the truth likely lies somewhere in between. Yet, there is little doubt that increasingly complex scientific issues have the potential to further tax the jury system, and that courts must seek new ways to help jurors deal with scientific evidence. To do so, courts will have to promote an active learning environment within the courtroom—in effect, turn courtrooms into classrooms. This new approach to jury trials is under way in some states today, pioneered by Arizona in its far-reaching 1995 jury rule changes including permitting jurors to ask questions, take notes, and in civil cases allowing jurors to discuss the evidence during the trial.3 Arizona's objective: improve the experience and decision making of jurors by redefining their role from passive observers to active participants, using applied, proven adult learning methods, and permitting information to unfold during the trial in more meaningful and understandable ways—in other words, to increase the potential of the "search for the truth."
As research on Arizona's jury reform experience progresses, there is growing evidence that the courtroom, turned juror-friendly classroom, is more conducive to juror comprehension and promotes ease in understanding complex concepts and data. If such is the case, must others wait for statewide system changes? The simple answer: no. Courts and lawyers already possess the means and discretion to enable juries to better carry out their vital roles. Judges and lawyers can independently recognize their roles as educators by embracing ground breaking jury reforms and introducing them in their own courts. These reforms will become increasingly important as genomic evidence appears ever more routinely in America's courtrooms.
Juries and complex cases
Over the past 30 to 40 years, the perceived performance of juries has been criticized, both in high-profile criminal cases and in complex civil litigation in antitrust, securities, intellectual property, and product liability cases. Critics have questioned whether a jury of untrained and inexperienced people can be a competent fact finder and decision maker in lengthy trials that require comprehension of substantial quantities of complex scientific, technical, or statistical evidence, and resolving the testimony of duplicative expert witnesses whose opinions conflict.
Moreover, it is alleged, juries in complex trials will have greater difficulty understanding and remembering the court's instructions, and properly applying the law to the facts. Faced with such a burden, say critics, jurors who are untrained in science and technology are ill-equipped for sound fact finding. As a result, critics allege, jurors will base their decisions less on the evidence and a careful consideration of the reliability of expert testimony, than on external cues, such as the perceived relative expertise and status of the expert witnesses, and will be more susceptible to "junk science" and emotional appeals.4
Intuitively then, we would expect juries to have enormous difficulties with the complex legal issues and scientific evidence that will confront the courts as disputes involving the strange, new world of human genetics and statistical probabilities become more commonplace. We would expect, as well, new proposals for replacing juries with such expert bodies as science courts and expert or "blue ribbon" panels. At the same time, however, a growing body of research on juries and their performance in both "simple" and complex cases is giving us a different picture.5 This research, based on case studies and "lab" or experimental studies, shows that jurors, rather than giving up in the face of voluminous evidence and conflicting expert opinions, take their fact-finding and decision-making responsibilities seriously.
The research shows that while certain elements of complex trials do tax jurors' comprehension and understanding, there is no firm evidence that their judgments have therefore been wrong. Jurors are in fact capable of resolving highly complex cases. These studies have also shown that factors such as length of trial, and evidentiary complexity in itself, are not necessarily the critical factors in jury performance in complex matters. The problem presented by conflicting testimony of experts hired by the respective parties, for example, is present in simple as well as complex cases. Finally, the research shows that jurors, rather than being passive participants in the trial process, are active decision makers and want to understand. Jurors actively process evidence, make inferences, use their common sense, have individual and common experiences that inform their decision making, and form opinions as a trial proceeds.6
What the research shows then, along with the experiments and experiences of active and concerned judges in complex cases, is that the trial process itself may be as much an impediment to jury comprehension and understanding as the complexity of the legal concepts and evidence, or the competencies of jurors.7 Many factors, including failure to follow instructions, confusing instructions, non-sequential presentation of evidence, "dueling" expert witnesses, evidentiary admissibility rulings, and attorney strategic errors, affect the jury's ability to follow and comprehend complex evidence. Researchers, and increasingly many progressive courts, suggest that reforming and improving the "decision making environment"8 can improve not only jury comprehension and performance, but juror satisfaction with their trial experience.
Challenging the current model
The Arizona Supreme Court's Committee on More Effective Use of Juries recognized these issues when it made 55 recommendations to reform the jury system, many of which resulted in the officially adopted comprehensive jury reform rules in 1995. In the introduction to Jurors: The Power of 12, its report to the supreme court, the Committee cited "unacceptably low levels of juror comprehension of the evidence" as one of the motivating factors in urging the Supreme Court to adopt its proposed jury reform rules.9 Arizona's reforms, designed to make jurors active participants during the trial, include juror note taking, pre-deliberation discussions of evidence during civil trials, and the right of jurors to ask written questions. The Arizona reforms also permit judges greater latitude in exercising their inherent powers to provide to each juror preliminary and final written jury instructions, as well as to open up a dialogue between the jurors, the judge, and the lawyers when a jury believes it is deadlocked or needs assistance. The result has been increased satisfaction with the judicial process by judges, lawyers, jurors, and litigants. For years, jury reforms such as note taking and question asking were opposed on the assumption that jurors would miss crucial pieces of evidence or assume the role of advocate rather than neutral fact-finder. The empirical evidence collected thus far, however, overwhelmingly indicates that such opportunities do not adversely affect the pace or outcome of trials.
It is intellectually arrogant for those in the legal system to assume that lay jurors are incapable of processing complex information. We have all been thrust into a technologically advanced world, and lawyers and judges are hardly better prepared for the task of sifting through scientific evidence than the jury. But common sense suggests that jury reform measures will aid understanding, and jurors themselves support reforms such as those described above.10 We should recognize that it makes little sense to oppose practices that make jurors more comfortable with complex scientific information. To drive the point home, we have often made the observation that it is difficult to imagine an academic setting in which taking notes and asking questions would not be permitted.
Fortunately, the tides are beginning to shift in the debate over jury reform. Already a number of states are adopting new rules; Arizona, Colorado, and California are just a few.11 In New York, much of the reform debate has centered on the selection, administration, and management of the jury, but substantive changes are not far behind. Reforms such as increased jury fees and security, and a juror hotline to report problems have been quite successful. However, the trend in these states and others is to expand beyond administrative concerns and attempt to improve jury deliberations and performance. These grassroots efforts led the American Bar Association in 1998 to adopt a number of jury reform ideals drafted by a Section of Litigation task force as part of its Civil Trial Practice Standards. In adopting these standards, the ABA recognized the need to provide juries, lawyers, and judges with the tools to increase jury comprehension in this era of increasingly complex evidentiary issues.
However, a complete overhaul of state and local jurisdictional rules is not necessary. These reforms can often be implemented, consistent with existing rules, at the discretion of the trial judge. Of course, when local rules conflict, those rules control, but most judges possess the inherent power to implement reforms in complex cases. For example, Rule 611 of the Federal (and Arizona) Rules of Evidence permit the judge to control the mode and order of questioning witnesses and presenting evidence. With the number of complex cases dramatically on the rise, judges and lawyers need to collaborate to help the jury become better fact-finders.
A practical guide
Many lawyers and judges seem to have forgotten the proper role of juries. Alexis de Tocqueville, the renowned historian, once said:
[t]he jury...may be regarded as a gratuitous public school, ever open, in which every juror learns his rights,...and becomes practically acquainted with the laws, which are brought within the reach of his capacity by the efforts of the bar, the advice of the judge, and even the passions of the parties...I look upon the [the jury] as one of the most efficacious means for the education of the people which society can employ.12
It is this idea of educating the jury, of treating the courtroom as a classroom, that judges and lawyers alike need to recapture. We urge all members of the legal profession to implement, on their own initiative, the appropriate reforms when cases require an understanding of complex scientific evidence.
Before we discuss individual reforms in more detail, it is important to note the role of judges in rigorously applying the rules of evidence. The judge plays a very important role in improving jury comprehension by appropriately screening evidence and admitting only that which meets the appropriate standards. The judge must scrupulously protect the jury from unreliable scientific evidence.13
Jury selection. Lawyers are often criticized for using their peremptory challenges to "dumb down" the jury. In complex cases, however, it is in the best interest of all concerned to select educated jurors and not strike persons based on the extent of their education. While there is little empirical evidence to demonstrate that more educated jurors are struck more often than less educated jurors, there does seem to be an unwritten rule of practice that professionals should be struck when possible. The authors themselves plead guilty to using that approach as trial lawyers.
Perhaps lawyers fear that highly educated individuals will dominate in the jury room and be able to persuade the jury to their side during deliberations. However, preliminary data suggest, and we believe, that jurors take their job seriously and will not be easily persuaded to a position with which they do not agree.14 Those lawyers who believe in "dumbing down" juries should adjust their views accordingly, and recognize the important role of jurors as fact finders and decision makers. Of course, both lawyers and judges must still attempt to detect jurors with prejudices or preconceived ideas, but they should also seek to empanel the best jurors available from the pool.
Juror note taking and notebooks. Of all the reforms discussed, allowing the jury to take notes during the trial must be the most common-sense and least controversial. Nevertheless many jurisdictions just don't get it. Research indicates that note taking does not distract jurors, nor does it create an undue influence on those jurors who choose not to take notes. Judges in Arizona instruct jurors that they are not obligated to take notes, and they tell the jury to pay attention to all aspects of the trial including witness demeanor and the documentary and testimonial evidence. The vast majority of courts recognize that it is within the sound discretion of the trial judge to permit jurors to take notes. Judges need to thoughtfully exercise their discretion and allow juror note taking in complex cases, and lawyers must urge judges to do so. Jurors need to be encouraged to take an active role in the trial. Allowing the jury to keep track of parties, witnesses, testimony, and evidence by taking notes will empower juries to improve their recall and understanding of all issues, simple and complex.
Jurors in complex cases should also be given a comprehensive notebook containing items such as simplified jury instructions, layouts of the courtroom with the names and locations of lawyers and parties, and glossaries of scientific terms or helpful scientific diagrams, photographs, charts, and background data of all types.
Better jury instructions. Judges historically instruct juries at the end of the trial. There are few rules or cases, however, that prohibit judges from instructing juries earlier. Judges in Arizona provide juries with pretrial instructions that, for example, define the elements of the alleged crime or define terms such as "negligence" and "fault." This permits the jury to understand the basic legal standards early in the case, refer to them during the trial, and then concentrate on the presentation of the evidence.
Jury instructions should be written in plain English. When drafting jury instructions, both judges and lawyers should avoid unnecessary legal jargon. In Arizona, the state bar's Civil Jury Instruction Committee even includes a linguistics professor from a local university. Jury instructions must also be tailored to the case at trial. Instead of using only pattern jury instructions, judges should work with counsel to draft case-specific instructions that include party names and actual facts in the case, without commenting on the evidence. Instructions should be given early in the case both orally and in writing for maximum comprehension and memory retention. The written instructions should be included in the jury notebook. Jurors need to understand the legal context of the evidence presented, and early instruction facilitates a better understanding of its legal relevance.
Finally, jurors should each be given a written copy of the final instructions and they should be allowed to have the instructions in the deliberation room. Arizona's rules require judges to provide each juror with a copy of all the jury instructions. After all, why should jurors have to pass a single copy when a few dollars can provide copies all around? And where is it written that jury instructions must only be oral?
Permitting the jury to ask written questions. When it comes to issues of scientific evidence, lawyers and judges collaborate to understand and narrow the issues before the court. They ask each other questions to clarify misunderstandings prior to trial, and will confer even during the trial. Yet, once the trial begins, jurors traditionally are not permitted to ask questions. It is time to end this nonsensical practice.
Jury questions should be written and given to court personnel before the witness leaves the courtroom. Counsel should be given the opportunity to object in a sidebar, or outside the hearing of the jury, and the jury should be instructed about the limitations on questions that can be asked. In Arizona, there have been no reports of problems with this type of procedure after thousands of trials over the last four years. A study reported in the March-April 1996 issue of Judicature found that jury questions helped jurors understand the facts and issues, that jurors did not ask inappropriate questions, and that jurors did not draw inappropriate inferences when their questions, due to counsel's objection, for example, were not asked.15
As the comments to the ABA Standards noted, state and federal courts have overwhelmingly recognized that it is within the sound discretion of the trial judge to allow juror questioning of witnesses. We encourage judges and lawyers to experiment with jury questions in complex cases. The empirical evidence, and our own experience, reveals that the fears and concerns about jury questions are unfounded. As two Arizona attorneys recently wrote, "Our experience [with juror questions] reinforces for us the effectiveness of juror questions in keeping the jury engaged and in improving the quality of our own trial presentations. The jurors' questions revealed areas of confusion or concern, enabling us to adjust our presentation accordingly."16
Juror discussion during civil trials. Perhaps one of the most controversial Arizona reforms at the time of its adoption, and still controversial today, is allowing jurors in civil cases to discuss the evidence prior to final deliberation. In Arizona, jurors are carefully instructed by the trial judge that they may discuss the case, so long as all members of the jury are present and they reserve judgment until final deliberations. The general consensus of the Arizona bench and bar is that this reform has been a success. In fact, the Committee on the More Effective Use of Jurors, in its second report to the Arizona Supreme Court (in June, 1998), recommended that the rules be expanded to allow pre-deliberation discussions during criminal trials. As of this writing, however, the supreme court has not adopted that recommendation.
Traditionally, the view has been that permitting jurors to discuss the evidence early in the trial will lead them to make up their minds before hearing both sides. Recent studies suggest that this is not true.17 In fact, some studies have gone so far as to say that requiring jurors to refrain from discussing evidence actually hinders their ability to process information.18 Pre-deliberation discussion can help improve juror comprehension, improve memory recall, and relieve the tension created by a forced atmosphere of silence with regard to the evidence presented at trial.19
Social scientists report that jurors naturally tend to actively process information as it is received. Therefore, it is not surprising to find that studies show that anywhere from 11 to 44 percent of jurors discuss the evidence among themselves during the trial despite judicial admonitions to avoid such discussion.20 Explicitly allowing pre-deliberation discussions, then, is really an acknowledgment of what often occurs naturally.
Perhaps surprising to some, Arizona's experience has shown that when one individual juror makes a preliminary judgment during pre-deliberation discussions, that judgment is often tested or challenged by the entire group.21 In United States v. Wexler (1987) Judge Ditter aptly explained that "jurors are concerned, responsible, conscientious citizens who take most seriously the job at hand." Like Judge Ditter, we believe the jurors are more interested in doing justice than in justifying their own loosely based preliminary conclusions, which are frequently subject to modification as a result of group discussions.
A recent study of jury discussions during Arizona trials found that jurors overwhelmingly support this reform and report that it has positive effects.22 Specifically, jurors said that discussions improved comprehension of evidence, that all jurors' views were considered, and evidence was remembered accurately. Additionally, only a very low percentage of participants in the study said that trial discussions encouraged jurors to make up their minds early on. The study also found that, among judges, lawyers, and jurors, support for this reform increases with experience. Permitting pre-deliberation discussion, more than any other reform, challenges the legal profession's traditional notions of jury behavior, but it is time to recognize the need for juries to have better tools in dealing with complex evidentiary issues.
Independent court appointed or stipulated experts. Unlike fingerprint or ballistic evidence, where it is easier to understand the samples juries are asked to compare, genetic evidence requires juries to sit through conflicting scientific interpretations from expert witnesses presented by the opposing parties. Early presentation of independent experts, either court appointed or stipulated, can help solve many of the problems presented by genetic evidence. Recent surveys suggest that judges favor appointing independent experts in complex cases. However, statistics show that the actual use of court appointed experts is relatively low.23 This situation is unfortunate because there are many advantages to be realized by the use of independent experts. For example, a case involving the admissibility of DNA evidence using a particular type of analysis was recently before the Arizona Superior Court. Both parties agreed to the appointment of a neutral court expert to testify about the procedures used in this analytial method. Substantial saving, in time and money, were realized by the appointment of the court expert. Judicial economy and fairness demand the use of innovative techniques in dealing with admittedly complex scientific issues.
In most jurisdictions trial judges have inherent authority to appoint experts as technical advisors to assist the court. In fact, judges may appoint expert witnesses for testimonial purposes under Rule 706 of the Federal Rules of Evidence and similar provisions in force in most states. However, the use of court appointed experts to serve as a jury tutor on the basics of, for example, DNA evidence, is an under-utilized tool.24 Pre-recorded video "lectures" may be another avenue to explore when considering how to educate jurors on issues of "common" scientific knowledge. The basic building blocks of DNA and the basic methods of DNA testing could be simplified and presented to the jury in such a fashion as to make it much less intimidating.25
Many lawyers may argue that "dueling experts" is the model courts should adhere to, based on the adversarial nature of our justice system. However, a recent study found that jurors do not rely on cross-examination of expert witnesses designed to point out flawed scientific methodology.26 The authors suggest that this is because jurors do not believe lawyers are sincere in their attempts to educate jurors, but rather see cross-examination as the lawyer's attempt to undermine the expert through any means possible.
Independent experts present an opportunity to not only improve juror comprehension and performance, but also decrease the substantial costs of expert witnesses, and increase judicial economy. The adversarial nature of the trial may be diminished, but that is actually a benefit, not a cost, according to independent experts considering jury reactions to lawyer cross-examination of opposing party witnesses. It is the judge's responsibility to be proactive in ensuring that the trial is a search for the truth, and that it is not about lawyers setting up roadblocks to that search.
Allow a dialogue between jurors, lawyers, and the judge during deliberations. In place of the traditional "pep talk" judges often give to deadlocked juries, Arizona explicitly provides for an opportunity for further instruction by the judge and argument by the parties. Why should the opportunity to educate jurors further stop once deliberations begin? Allowing additional evidence, argument by counsel, or providing further instruction is not problematic, legally or pragmatically. Of course, judges must be careful not to influence jurors and need to limit further inquiries only to those issues that confuse or divide the jury. Once again, there are many cases approving the judge's inherent authority to reopen a case for additional evidence or argument where the jury needs further admissible evidence to reach a verdict, or to determine if a deadlock is unavoidable.27
Opening the courtroom to more creative learning. Increasingly, the Human Genome Project's Ethical, Legal and Social Implications Program is sensitizing the judicial and legal community about the changing rule of the law in light of new genetic discoveries and testing methods. Primers reviewing DNA and genome science have been written, memorable cartoon drawings simplify sophisticated concepts,28 and video background resources explaining genetics in meaningful non-scientific ways are growing in number.
Further, difficult concepts can be reduced to plain English and conveyed to juries through innovative technologies, including live, videotaped, or interactive Internet-based testimony. These approaches can easily be presented while simultaneously ensuring that complex scientific evidence is afforded the utmost of seriousness.
Educating the jury early in the trial, by using court appointed experts, better written jury instructions, jury notebooks, and basic adult education techniques, will provide a foundation for later testimony of experts presented by the lawyers. Jurors who have been tutored early about complex scientific issues will be in a better position to judge both the content and character of dueling experts.
Two central participants in the courtroom are the ultimate beneficiaries of reform-oriented jury approaches when heavy doses of scientific evidence are the subject of an unfolding courtroom drama: jurors, and more importantly, litigants. Contemporary behavioral research, and Arizona's jury reform experience, substantiate that comprehension and understanding are significantly enhanced when information is actively processed. Most courts already possess the tools to implement the educational techniques discussed above. Whether through system-wide jury reform or the efforts of individual trial judges and trial lawyers, a more jury-centered trial will not only allow jurors to actively and intelligently participate in the fact-finding and decision-making process, but also give the litigants a better truth-finding forum.
Robert D. Myers is Presiding Judge of the Arizona Superior Court in Maricopa County.
Ronald S. Reinstein is Associate Presiding Judge of the Arizona Superior Court in Maricopa County.
Gordon M. Griller is court administrator, Arizona Superior Court in Maricopa County and a member of the Board of Directors of the American Judicature Society.
The authors wish to thank Timothy D. Keller, a law researcher for Judge Robert D. Myers, and Richard Teenstra, assistant director of the Maricopa County Superior Court Law Library, for their assistance.
1. Department of Energy, Office of Biological and Environmental Research, Life Sciences Division, Human Genome Research: An Introduction (visited Sept. 2, 1999) http://www.science.doe.gov/.
2. Denno, Legal Implications of Genetics and Crime Research, in Bock and Goode, eds., Genetics of Criminal and Antisocial Behaviour 235 (Chichester, N.Y.: Wiley, 1996).
3. See Arizona Supreme Court Orders, Nos. R-94-0031, R-92-004 (1995).
4. See Adler, The Jury: Trial and Error in the American Courtroom (New York: Times Books, 1994); Jury Comprehension in Complex Cases: Report of a Special Committee of the ABA Litgation Section (Chicago: American Bar Association, 1989).
5. For a review of criticisms of civil jury competencies and the jury research literature, see Lempert, Civil Juries and Complex Cases: Taking Stock after Twelve Years, in Litan, ed., Verdict: Assessing the Civil Jury System 181-247 (Washington, D.C.: Brookings Institution, 1993); Vidmar, The Performance of the American Civil Jury: An Empirical Perspective, 40 Ariz. L. Rev. 849 (1998); Cecil, Hans and Wiggins, Citizen Comprehension of Difficult Issues: Lessons from Civil Jury Trials, 40 Am. U. L. Rev. 727 (1991).
6. Hans, Hannaford and Munsterman, The Arizona Jury Reform Permitting Civil Jury Trial Discussions: The Views of Trial Participants, Judges, and Jurors, 32 U. Mich. J.L. Reform 349 (1999).
7. See Dann, "Learning Lessons" and "Speaking Rights": Creating Educated and Democratic Juries, 68 Ind. L.J. 1229 (1993).
8. Cecil, Hans and Wiggins, supra n. 5, at 765.
9. Jurors: The Power of 12, Report of the Arizona Supreme Court Committee On More Effective Use of Juries (November 1994).
10. Hans, Hannaford and Munsterman, supra n. 6, at 371-372.
11. For a review of state jury reform efforts, see Munsterman, A brief history of state jury reform efforts, 79 Judicature 216 (1996); Murphy, et al, Managing Notorious Trials (Williamsburg, Va.: National Center for State Courts, 1998); Enhancing the Jury System: A Guidebook for Jury Reform (Chicago: American Judicature Society, 1999).
12. de Tocqueville, Democracy in America 295-296 (Vintage ed. 1945).
13. Daubert v. Merrell Dow Pharm. Inc., 509 U. S. 579 (1993).
14. Hans, Hannaford and Munsterman, supra n. 6.
15. Heuer and Penrod, Increasing juror participation in trials through note taking and question asking, 79 Judicature 256, 260-261 (1996).
16. Cabot and Coleman, Arizona's 1995 Jury Reform Can be Deemed a Success, Arizona Journal, July 12, 1999, at 6.
17. See Hans, Hannaford and Munsterman, supra n. 6; Hannaford, Hans and Munsterman, "Permitting Jury Discussions During Trial: Impact of the Arizona Reform" 9 (1998) (unpublished manuscript, on file with the authors).
18. Chilton and Henley, Improving the Jury System, Jury Instructions: Helping Jurors Understand the Evidence and the Law, §II, PLRI Reports (Spring 1996) http://www.uchastings.edu/plri/spr96tex/juryinst.html.
19. Hans, Hannaford and Munsterman, supra n. 6; Hannaford, Hans and Munsterman, supra n. 17; Chilton and Henley, supra n. 18.
20. Chilton and Henley, supra n. 18.
21. Myers and Griller, Educating Jurors Means Better Trials: Jury Reform in Arizona, 36 Judges J. 13-17, 51 (Fall 1997).
22. Hans, Hannaford and Munsterman, supra n. 6.
23. Sanders, Scientifically Complex Cases, Trial by Jury, and the Erosion of Adversarial Processes, 48 DePaul L. Rev. 355, 378-379 (1998).
24. The Evaluation of Forensic DNA Evidence 169-171 (Washington, D.C.: National Research Council, 1996).
25. For examples of excellent illustrations and explanations, see Hoagland and Dotson, The Way Life Works (New York: Time Books, 1995).
26. Kovera, McAuliff and Hebert, Reasoning About Scientific Evidence: Effects of Juror Gender and Evidence Quality on Juror Decisions in a Hostile Work Environment Case, 84 J. of Applied Psychology 362, 372-373 (1999).
27. Myers and Griller, supra n. 21, at 16-17.
28. See Hoagland and Dotson, supra n. 25.
|The online presentation of this publication is a special feature of the Human Genome Project Information Web site.| |
By Gerald Forbes
While the South Burbank pool is being cited as an example of the economic advantages of unitization,1 it is an interesting fact that the first two decades of oil development in the Osage Reservation contained similar features of community interest and concerted control. Between the years 1896 and 1916 the petroleum of the Osage Reservation was developed under a single contract, known as the Foster Lease, the only enduring and successful "blanket" lease in Mid-Continent's history. The Foster Lease undoubtedly was a monopoly,2 but just as certainly it was an instrument of conservation.
The former Osage Reservation, the present Oklahoma county by that name, contains about one and a half million acres that were bought from the Cherokees, preparatory to moving the Osages from Kansas in 1872. This territory is bounded by the ninety-sixth meridian on the east, the Arkansas River and the former Creek Nation on the south and west, and Kansas on the north.3
It was nearly twenty years after the Osages had bought the land that the possibility of producing petroleum began to be investigated. The American Civil War had interrupted the beginning of the oil industry in Kansas, but in the final decade of the century the production of petroleum became an established industry in that state. In 1895 the oil production of Kansas was 44,430 barrels.4 Among those persons interested, in the oil industry was Henry Foster, who had moved to Independence, Kansas, from Rhode Island. Foster suspected the presence of petroleum beneath the land of the Osages. In 1895 he applied to the Secretary of the De-
1John J. Arthur, "Unitization vs. Competition," The Oil Weekly, V. 83, No. 2, September 21, 1936, pp. 22-26; The Oil Weekly, V. 82, No. 13, September 7, 1936, p. 51.
2Kate P. Burwell, "Richest People in the World," Sturm's Oklahoma Magazine, II, No. 4, pp. 89-93; United States Geological Survey, Mineral Resources of the United States, 1905, p. 885.
4United States Geological Survey, Mineral Resources of the United States, 1889-90, 355; University Geological Survey of Kansas, IX, Special Report on Oil and Gas, 1908, pp. 21-23.
partment of the Interior for a lease of the Osage Reservation. Henry Foster died before the contract was consummated, but his brother, Edwin B. Foster, assumed his interests and obligations. The lease finally was signed, March 16, 1896, by Edwin B. Foster and the Osage National Council, with James Bigheart, principal chief of the Osages, Saucy Chief, president of the Council, and several other Indians writing their names or making their "X's."5
The terms of the contract, which soon became known as a "blanket lease," conferred on Foster the exclusive right of producing oil in the entire Osage Reservation. The term was for ten years. Foster, in turn, agreed to pay a royalty of ten per cent of all crude petroleum removed from the ground and fifty dollars a year for each gas well, as long as it was used. The royalty was to be based on the market value at the place of production, and was to be paid to the National Treasurer of the Osage Nation. Foster further agreed to settle the royalty accounts between the fifth and tenth days of January, April, July, and October.6 Even at this early period there was dissention, and in less than a month a protest was filed with the Secretary of the Interior. The leading protestant was Saucy Chief, who had placed his "X" on the original Foster Lease. The protest was not attested and not clearly genuine. It declared that a full council had not been present when the leasing had been discussed and that the contract did not represent the wishes of a majority of the Osage Tribe. An investigation followed, and the Osage Agent reported that two white men had taken about fifty Indians across the Arkansas River to Cleveland, Oklahoma Territory, where the Osages had been induced with whiskey to sign a protest to the Foster Lease.7
Edwin B. Foster and the heirs of Henry Foster, having organized the Phoenix Oil Company, arranged with McBride and
5Osage Indian Archives, Pawhuska, Oklahoma, D. M. Browning to Henry Foster, January 24, 1896; Hines, E. P., Osage County, in Snider, L. C., Oil and Gas in the Mid-Continent Field, p. 208; Kappler, Charles J., Indian Affairs, Laws and Treaties, III, 1913, p. 137.
6Exact copy of the original lease—Mining Lease, Osage Agency, Oklahoma Territory, 1896, for Prospecting and Mining for Oil and Gas upon the Osage Reservation, Oklahoma Territory.
7Osage Indian Archives, D. M. Browning to U. S. Indian Inspector Duncan, April 6, 1896; D. M. Browning to Acting Agent Freeman, June 13, 1896.
Bloom, drillers of Independence, Kansas, to put down a well three or four miles south of Chautauqua Springs, Kansas, likely near the present town of Boulanger, Oklahoma.8 This well was shallow but it produced about fifty barrels of oil daily, not enough at that time to be commercially valuable, so it was capped. The rumor that the well was an excellent one became current. It was rumored that this first Osage well had been closed to permit the owners to acquire leases cheaply in the Oklahoma Territory. The first Osage well was drilled in 1897. It was in 1899 that the Osage Oil Company, another Foster concern, drilled on the eastern side of the Osage Reservation near Bartlesville, a town in the Cherokee Nation. The first well of the Osage Oil Company showed prospects of petroleum and the second well was a good producer. Several dry or nearly dry holes were drilled, but the seventh well of the group was the best producer in the entire Kansas-Indian Territory oil field.9
By 1900 Foster had done little to develop the petroleum industry in the Osage Reservation, but in that year arrangements were completed for subleasing the land in large blocks. The entire reservation was divided into tracts half a mile wide and three miles east to west. These rectangles were numbered consecutively and those in the eastern part of the Reservation were offered to sublessees on a bonus and royalty basis. The sublessees were required to pay the Foster interests a one-eighth (later one-sixth) royalty and a bonus of one to five dollars an acre.10 The next year the Foster interests were consolidated in the Indian Territory Illuminating Oil Company (usually called the I. T. I. O.) which was incorporated at Trenton, New Jersey, with a capitalization of three million dollars. This new company was authorized to own and control all the rights and properties of the Osage and Phoenix Oil Companies.11 It was the I. T. I. O. that handled the subleasing
8Hines, loc. cit. p. 208; Hutchison, L.L., Preliminary Report on Rock Asphalt, Asphalite, Petroleum and Natural Gas in Oklahoma, Oklahoma Geological Survey, Bulletin No. 2, 1911, p. 167; Tidal Topics, III, Tidal Oil Company, 1919, p. 15.
11The Tulsa Democrat, Tulsa, Indian Territory, December 27, 1901, The Osage Journal, January 2, 1902.
of the Osage Reservation, and buyers of drilling rights were sought in New York. The first well drilled by a sublessee was financed by the Almeda Oil Company on Lot 40. The Indian Territory Illuminating Oil Company announced that it planned to drill wells itself at the rate of one every twenty days, that the Standard Oil Company would buy the crude oil production at its refinery at Neodesha, Kansas, for eighty-eight cents a barrel, and that leases had been sold to New York and St. Louis companies covering rights on about six thousand acres of land. The I. T. I. O. further called attention to the quality of the crude oil which caused it to yield a high percentage of kerosene. (The name of the company itself calls attention to the fact that gasoline then was not of first importance.) The average depth of the wells was thirteen hundred feet, which made them relatively inexpensive to drill.12
Drilling in the Osage Reservation was comparatively rapid after the system of subleasing had been perfected. During 1902 the rail shipments of crude oil to the Neodesha refinery amounted to 37,000 barrels, which was the production of thirteen wells, six of which had been drilled in 1902. By January, 1903, thirty wells had been completed by the I. T. I. O. and its sublessees. Seventeen of the thirty wells produced oil, two gas, and eleven were dry holes. A year later 361 wells had been completed, and 243 were producing oil, twenty-one gas, and ninety-seven were dry. By the beginning of 1906 there had been 783 wells drilled—544 producing oil, forty-one gas, and 198 were dry. The oil production was: 1903—56,905 barrels; 1904—652,479 barrels; 1905—3,421,478 barrels; 1906—5,219,106 barrels. The average daily production of the Osage wells in 1905 was about 15,000 barrels. In 1905 there were 687,000 acres of the Osage Reservation under the control of the sublessees.13 That year the I. T. I. O. announced that it had disbursed $2,686,627 in connection with the "blanket lease."
By the terms of the contract with Foster, the Osages were to receive one-tenth royalty (later changed to one-eighth) while the I. T. I. O. Company required one-eighth (later changed to one-
13United States Geological Survey, Mineral Resources of the United States, 1905, p. 855; 1906, p. 858; 1914, pp. 1009-1010; Tidal Topics, III, p. 15.
sixth) royalty of its sublessees, making a profit of one-fortieth (later one-twenty-fourth) in addition to rentals and bonuses. There were less than twenty-five hundred members of the Osage tribe on the official rolls. The rolls contained Indians on the list January 1, 1906, and all children born to them by July, 1907, and those children of white fathers who had not been enrolled previously. There was no distinction between males and females, age, or degree of Indian blood. The equal share which each member of the tribe received from the communal mineral receipts was known as a headright. Headrights, it was provided by law, could be inherited, subdivided or consolidated, and as time passed different members of the tribe did not receive equal shares, as was the case at the time of the Osage allotment. This allotment differed from that of the other Oklahoma tribes, for it provided that only the surface of the land be held in severalty while the minerals of the subsurface remained communal property. As the sublessees of the I. T. I. O. developed the oil industry, the royalties of the Osage tribe mounted and were divided into headright payments.14
The days of the quarterly payments at Pawhuska, seat of the Osage agency, were colorful. On the first and second days of the payments, the full bloods received their monies; then the mixed-bloods were paid on the following two or three days. By 1906 the quarterly payment period kept force of eight men busy for four or five days. The merchants and professional men of Pawhuska who had extended credit to the Indians were on hand to collect their bills before the Osages had spent their money elsewhere. The amount of the payment depended on the number of barrels of oil taken from the ground, the number of gas wells being used, and the market price of petroleum. Accurate figures on the receipts from oil and gas are difficult to acquire, for the Osages also received payments for grazing permits, pipe line damages, and other revenues. Between July 1, 1904, and May 13, 1905, a total of $108,567 was paid to the Osages as oil and gas royalties.15
14United States Statutes At Large, XXXIV, p. 540; Kappler, op. cit., p. 256; United States Geological Survey, Mineral Resources of the United States, 1906, p. 855; Daniel, L. H. "The Osage Nation," The Texaco Star, V. Nos. 7-8, pp. 10-14.
Congress began considering the renewal of the Foster Lease in 1905, although it did not expire until March, 1906. Several of the tribal leaders went to Washington to watch the action of Congress, and there were some who wished to prevent renewal of the contract.16 There were oil operators who called attention to the profits they believed the I. T. I. O. company was making and objected that one firm should have such a monopoly. After an investigation, Congress compromised by renewing the Foster Lease and all the subleases made by the I. T. I. O. on a total area of six hundred and eighty thousand acres on the eastern side of the Reservation. All the original conditions of the Foster Lease were to apply for another decade, with the exception that gas well royalty was increased from fifty to one hundred dollars for each well. The status of the western half of the Reservation was left undetermined until 1912. The renewal with reduced acreage left the Indian Territory Illuminating Oil Company with only 2,060 acres that had not been subleased, and caused that firm to lease from its own sublessees.17
Before 1904 the Osage oil was transported by railroad, but in that year the Department of the Interior approved two applications for pipelines to move the crude petroleum. The amount of damages to be paid the Osages puzzled the Federal officials, for there was no precedent for laying pipelines across Indian lands. The Prairie Oil and Gas Company wanted to lay a line to the refinery at Neodesha, while Guffey and Galey sought to pipe gas to Tulsa.18 Damages were fixed at ten cents a rod. In 1905 the Prairie constructed the "Cleveland discharge" line, which connected the Osage wells near Cleveland, Oklahoma Territory, with the trunk line to Kansas. Another outlet for the Osage petroleum appeared with the construction of a refinery by the Uncle Sam Oil Company at Cherryvale, Kansas. The disagreements of the Uncle Sam and the
17"History of twenty-three Years of Oil and Gas Development in the Osage," "National Petroleum News," V. No. 11, pp. 66-68; Kappler, op. cit., p. 137; Osage Indian Archives, Memorandum, p. 1; Osage Indian Archives, C. F. Larrabee to Frank Frantz, June 7, 1905; Hines, E.P., loc. cit., p. 208; Department of the Interior, Commissioner of Indian Affairs, Annual Report, I, p. 307.
18The Osage Journal, March 18, 1905; The Cherokee Advocate, Tahlequah, Indian Territory, April 4, 1903.
Standard companies were dragged through the courts for years.19 In 1910 the Gulf Pipe Line Company became a buyer of Osage oil, since inadequate transportation facilities had resulted in 1909 in a decrease of production.20 Drilling and production received no more setbacks until 1915, when little drilling was done because of the uncertainty resulting from the struggle over the second renewal of the Foster Lease.
The disposition of the mineral rights in the western half of the Osage county (Oklahoma became a state in 1907) became a pressing question in 1911. A committee of Osages urged the National Council to lease the western land on terms that would be more profitable to the Indians. Royalties of one-third and one-sixth were suggested. Since the Osages were interested in farming and ranching, as well as oil, it was argued that no company should be permitted to drill for oil without the "written consent" of the allottee on whose land the well was desired. After revising some of the suggestions of the committee, the Osage National Council went on record as favoring sealed bids for leases. Sealed bids would prevent leasing except at specific times, and then the lease would go to the highest bidder.21
While this discussion was current among the Indians, some oil operators met at Tulsa and decided on a plan for leasing the western side of Osage County. They proposed the organization of a large company of independent operators, each of whom would be on an equal cooperative footing. Such a company, the oil men believed, would be financially able to contract for the entire unleased acreage of Osage lands. They believed this company could deal pleasantly with the Department of the Interior. The financing of this huge company was expected to be comparatively simple, and it was argued that such a concern would be able to dictate favorable terms to crude oil buyers and thereby gain a profitable
19Osage Indian Archives, C.F. Larrabee to Frank Frantz, January 12, and January 14, 1905; The Muskogee Times Democrat, Muskogee, Indian Territory, January 8, 1907.
21The Osage Journal, January 5, May 25, August 17, September 14, and October 19, 1911; Senate Document 487, 62 Congress, 2 Sess.
price for the petroleum. Among the leaders of this plan were P. J. White, Harry Sinclair, E. R. Kemp, and David Gunsberg.22
Samuel Adams, Assistant Secretary of the Department of the Interior asked those who were interested in leasing Osage land to communicate with him.23 The Osage National Council went to Washington to confer with Adams. The proponents of the giant organization of independent producers sent representatives. Many oil men favored neither the plan of the independents nor that of the Osage Council, so a mass meeting was called at Tulsa to protest the organization of the giant cooperative firm. Some believed that the Osage oil long had been a menace to the price of petroleum, and they did not look kindly on any plan to further the production. They suggested that a plan be adopted to discover whether any oil existed in the western side of the Osage County. Several operators believed that all the oil of the Osages had been discovered. Another group, led by E. W. Marland and F. A. Gillespie, opposed any plan involving one big lease. They favored leasing the western side of the county in blocks as small as 160 acres.24
In May, 1912, the Osage National Council directed the principal chief to sign four leases that would cover virtually the entire western part of the county. In these leases were several ideas which the Osages desired, including the "written consent clause," the maintenance by the leasing companies of offices at Pawhuska, and the retention in the county of all the gas. (It was believed that the retention of the gas in the county would induce industries to come.) The leases were issued to four men, one of whom was H. H. Tucker of the Uncle Sam Oil Company, who was reported to be an adopted member of the Osage tribe. The Secretary of the Department of the Interior refused to accept these leases because no provision was made for the supervision of the Federal government. He also frowned on the "written consent clause."25
Despite the fact that Assistant Secretary Adams had said that he would not recognize their election, Bacon Rind, as Principal
25The Osage Journal, March 14 and May 23, 1912; The Tulsa World, March 16, May 25 and June 19, 1912.
Chief, and Red Eagle, as Assistant Chief, celebrated their election in July of that year (1912).26 Under the guidance of Bacon Rind and Red Eagle, the Osage National Council joined the Uncle Sam Oil Company in publicly presenting a petition to President Taft asking that the entire unleased portion of the Osage lands be leased to Tucker's company. The Department of the Interior concluded the opposition among the Indians by promptly removing from office both Bacon Rind and Red Eagle, as well as the entire National Council. Tucker responded with a final threat to President Taft that the twelve thousand stockholders of the Uncle Sam Oil Company would remember the refusal of the president to override the decision of the Department of the Interior. He vowed that the stockholders would use their influence to prevent Taft's reelection in November.27
The final decision of the Department of the Interior, issued July 13, 1912, involved elements of several of the plans suggested for the disposal of the west side mineral rights. The land was to be leased in tracts varying from three hundred to 5,120 acres, but no person was to have more than 25,000 acres. The United States Agent at Pawhuska was required periodically to advertise specific tracts for leasing on sealed bids. A person wishing to lease a tract was required to request in writing that the land be offered for bidding. Each bid was to be accompanied by a certified check for ten per cent of the bonus and the first year's rental. All leases were to endure for ten years from the date of approval by the Department of the Interior, providing no lease extended beyond April 8, 1931. The royalty on gas was fixed at one-sixth of the market value at the well, while on petroleum it was set at one-sixth of the gross production at the actual market value. Heretofore the royalty on oil had been one-eighth. Oil men who had been paying one-eighth royalty on oil produced on the land of the Five Tribes objected to giving one-sixth to the Osages, but that was the share which the I. T. I. O. had been receiving from its sublessees. A compromise was reached on the "written consent clause" whereby cultivated lands and homesteads were protected from oil prospec-
tors. Producers strongly condemned the new regulations and the Osage National Council.28
The conflict over leasing the west side of the county hardly had ended before it was time for the renewal of the Foster Lease on the east side of the Osage Reservation. The I. T. I. O. minimized the profit it received from the Foster Lease, but in June, 1914, a renewal of the lease was asked. The request of the I. T. I. O. was supported by the company's sublessees. The next month the Osage National Council requested that no blanket lease be approved for the land then held by the I. T. I. O. The leasing company issued a financial statement to show the benefits that it had brought the Osages. The statement said that the I. T. I. O. had received over two million dollars in seventeen years, but that more than a million dollars had been paid to the Indians. The company cited the fact that it had furnished more than one hundred thousand dollars worth of gas free to operating companies. The statement of the I. T. I. O. indicated that the company had spent more in developing the Osage petroleum than it had received from the sublessees.29
When 1915 opened it was clear that some decision must be made regarding the Foster Lease. Secretary Lane of the Department of the Interior called a public hearing at Washington to discuss the lease. Members of the Osage National Council, officials of the Indian administration, and oil operators attended.30 Charles N. Haskell, first governor of Oklahoma, appeared for P. J. White and Harry Sinclair, and declared that the decision would affect the entire Mid-Continent. He asserted that the I. T. I. O. would develop the oil industry in an orderly manner, but that if the district were thrown open to competitive drilling there would be a
28Oklahoma Geological Survey, Bulletin No. 19, Part 1, Petroleum and Natural Gas in Oklahoma, 1915, p. 32; Department of the Interior, Regulations to Govern the Leasing of Lands in the Osage Reservation, Oklahoma, for Oil Gas, and Mining Purposes, 1912, pp. 1-4.
29Estimate of Profit and Loss under the Leases and Subleases of the Indian Territory Illuminating Oil Company in the Osage Reservation, compiled by Charles F. Leech, (nd) Osage Indian Archives.
30The Oil and Gas Journal, February 11, 1915, p. 2; Osage Indian Archives, Cato Sells to J. George Wright, February 10, 1915, The Osage Journal, February 11, 1915.
flood of oil that would swamp the marketing facilities. Charles Owen, in a letter that was made a part of the record, took the stand that if the I. T. I. O. were to be protected for the pioneer development, its lease should be renewed where it actually had put down wells, not subleased the land to other companies. Some sublessees objected to the policy of the I. T. I. O. in separating the oil and gas rights, for they argued that they had found the gas, but now that a market was available the I. T. I. O. held it. (By the Foster Lease, the I. T. I. O. owned all gas discovered.) The Osage National Council demanded that leases be made directly with the operating companies without the I. T. I. O. as an intermediary.31 The hearing was concluded in June and the Department of the Interior refused to renew the Foster Lease, deciding to eliminate the I. T. I. O. except as a producing company.
The new regulations provided that the east side of the county be broken up into quarter-section units combined in such a way that none would exceed an aggregate of 4,800 acres, except in such units where producing wells were capable of averaging twenty-five barrels a day on July 1, 1915. These units were to be offered at public auction for lease by the Osages under the supervision of the Department of the Interior. Oil and gas rights still were to be kept separately. The royalty on oil was fixed at one-sixth, except on quarter-sections where the average daily production equalled or exceeded one hundred barrels daily. There the royalty was one-fifth. Former sublessees of the I. T. I. O. were allowed to keep those quarter-sections they then were developing provided there would not be a total exceeding 4,800 acres.32 In general these rules were much the same as those governing the oil leases in the lands of the Five Civilized Tribes.
March 16, 1916, the Foster Lease expired, ending the only successful blanket lease of the lands of an Indian tribe. The lease
31Osage Indian Archives, J. George Wright to Cato Sells, March 2, 1915, Stenographer's Minutes of Hearing Before Cato Sells, Commissioner of Indian Affairs, in the Matter of the so-called Foster Lease on Oil and Gas Property Owned by the Osage Indians of Oklahoma, Washington, March, 1915, pp. 675-680, 682-683, 686-687, 692-694-703, 710-711, 715, 730; The Osage Journal, May 15, 1915.
was a monopoly, but it had good features for all concerned. The Osage lands continued to produce an increasing volume of oil until 1923, whereas the immense deposits in the Creek Nation (Glenn Pool, Cushing Pool, Okmulgee County) were dissipated very rapidly. In the Osage area there was a tendency to avoid competitive drilling because of the large leases. The Osage gas was conserved, thereby retaining much of the natural pressure. Gross overproduction never was one of the evils found in the district. The Osages themselves certainly benefited under the Foster Lease, although their individual wealth generally was over-estimated. |
Significance and Use
Sediment provides habitat for many aquatic organisms and is a major repository for many of the more persistent chemicals that are introduced into surface waters. In the aquatic environment, most anthropogenic chemicals and waste materials including toxic organic and inorganic chemicals eventually accumulate in sediment. Mounting evidences exists of environmental degradation in areas where USEPA Water Quality Criteria (WQC; Stephan et al.(67)) are not exceeded, yet organisms in or near sediments are adversely affected Chapman, 1989 (68). The WQC were developed to protect organisms in the water column and were not directed toward protecting organisms in sediment. Concentrations of contaminants in sediment may be several orders of magnitude higher than in the overlying water; however, whole sediment concentrations have not been strongly correlated to bioavailability Burton, 1991(69). Partitioning or sorption of a compound between water and sediment may depend on many factors including: aqueous solubility, pH, redox, affinity for sediment organic carbon and dissolved organic carbon, grain size of the sediment, sediment mineral constituents (oxides of iron, manganese, and aluminum), and the quantity of acid volatile sulfides in sediment Di Toro et al. 1991(70) Giesy et al. 1988 (71). Although certain chemicals are highly sorbed to sediment, these compounds may still be available to the biota. Chemicals in sediments may be directly toxic to aquatic life or can be a source of chemicals for bioaccumulation in the food chain.
The objective of a sediment test is to determine whether chemicals in sediment are harmful to or are bioaccumulated by benthic organisms. The tests can be used to measure interactive toxic effects of complex chemical mixtures in sediment. Furthermore, knowledge of specific pathways of interactions among sediments and test organisms is not necessary to conduct the tests Kemp et al. 1988, (72). Sediment tests can be used to: (1) determine the relationship between toxic effects and bioavailability, (2) investigate interactions among chemicals, (3) compare the sensitivities of different organisms, (4) determine spatial and temporal distribution of contamination, (5) evaluate hazards of dredged material, (6) measure toxicity as part of product licensing or safety testing, (7) rank areas for clean up, and (8) estimate the effectiveness of remediation or management practices.
A variety of methods have been developed for assessing the toxicity of chemicals in sediments using amphipods, midges, polychaetes, oligochaetes, mayflies, or cladocerans (Test Method E 1706, Guide E 1525, Guide E 1850; Annex A1, Annex A2; USEPA, 2000 (73), EPA 1994b, (74), Environment Canada 1997a, (75), Enviroment Canada 1997b,(76)). Several endpoints are suggested in these methods to measure potential effects of contaminants in sediment including survival, growth, behavior, or reproduction; however, survival of test organisms in 10-day exposures is the endpoint most commonly reported. These short-term exposures that only measure effects on survival can be used to identify high levels of contamination in sediments, but may not be able to identify moderate levels of contamination in sediments (USEPA USEPA, 2000 (73); Sibley et al.1996, (77); Sibley et al.1997a, (78); Sibley et al.1997b, (79); Benoit et al.1997, (80); Ingersoll et al.1998, (81)). Sublethal endpoints in sediment tests might also prove to be better estimates of responses of benthic communities to contaminants in the field, Kembel et al. 1994 (82). Insufficient information is available to determine if the long-term test conducted with Leptocheirus plumulosus (Annex A2) is more sensitive than 10-d toxicity tests conducted with this or other species.
The decision to conduct short-term or long-term toxicity tests depends on the goal of the assessment. In some instances, sufficient information may be gained by measuring sublethal endpoints in 10-day tests. In other instances, the 10-day tests could be used to screen samples for toxicity before long-term tests are conducted. While the long-term tests are needed to determine direct effects on reproduction, measurement of growth in these toxicity tests may serve as an indirect estimate of reproductive effects of contaminants associated with sediments (Annex A1).
Use of sublethal endpoints for assessment of contaminant risk is not unique to toxicity testing with sediments. Numerous regulatory programs require the use of sublethal endpoints in the decision-making process (Pittinger and Adams, 1997, (83)) including: (1) Water Quality Criteria (and State Standards); (2) National Pollution Discharge Elimination System (NPDES) effluent monitoring (including chemical-specific limits and sublethal endpoints in toxicity tests); (3) Federal Insecticide, Rodenticide and Fungicide Act (FIFRA) and the Toxic Substances Control Act (TSCA, tiered assessment includes several sublethal endpoints with fish and aquatic invertebrates); (4) Superfund (Comprehensive Environmental Responses, Compensation and Liability Act; CERCLA); (5) Organization of Economic Cooperation and Development (OECD, sublethal toxicity testing with fish and invertebrates); (6) European Economic Community (EC, sublethal toxicity testing with fish and invertebrates); and (7) the Paris Commission (behavioral endpoints).
Results of toxicity tests on sediments spiked at different concentrations of chemicals can be used to establish cause and effect relationships between chemicals and biological responses. Results of toxicity tests with test materials spiked into sediments at different concentrations may be reported in terms of an LC50 (median lethal concentration), an EC50 (median effect concentration), an IC50 (inhibition concentration), or as a NOEC (no observed effect concentration) or LOEC (lowest observed effect concentration). However, spiked sediment may not be representative of chemicals associated with sediment in the field. Mixing time Stemmer et al. 1990b, (84), aging ( Landrum et al. 1989,(85), Word et al. 1987, (86), Landrum et al., 1992,(87)), and the chemical form of the material can affect responses of test organisms in spiked sediment tests.
Evaluating effect concentrations for chemicals in sediment requires knowledge of factors controlling their bioavailability. Similar concentrations of a chemical in units of mass of chemical per mass of sediment dry weight often exhibit a range in toxicity in different sediments Di Toro et al. 1990, (88) Di Toro et al. 1991,(70). Effect concentrations of chemicals in sediment have been correlated to interstitial water concentrations, and effect concentrations in interstitial water are often similar to effect concentrations in water-only exposures. The bioavailability of nonionic organic compounds in sediment is often inversely correlated with the organic carbon concentration. Whatever the route of exposure, these correlations of effect concentrations to interstitial water concentrations indicate that predicted or measured concentrations in interstitial water can be used to quantify the exposure concentration to an organism. Therefore, information on partitioning of chemicals between solid and liquid phases of sediment is useful for establishing effect concentrations Di Toro et al. 1991, (70).
Field surveys can be designed to provide either a qualitative reconnaissance of the distribution of sediment contamination or a quantitative statistical comparison of contamination among sites.
Surveys of sediment toxicity are usually part of more comprehensive analyses of biological, chemical, geological, and hydrographic data. Statistical correlations may be improved and sampling costs may be reduced if subsamples are taken simultaneously for sediment tests, chemical analyses, and benthic community structure.
Table 2 lists several approaches the USEPA has considered for the assessment of sediment quality USEPA, 1992, (89). These approaches include: (1) equilibrium partitioning, (2) tissue residues, (3) interstitial water toxicity, (4) whole-sediment toxicity and sediment-spiking tests, (5) benthic community structure, (6) effect ranges (for example, effect range median, ERM), and (7) sediment quality triad (see USEPA, 1989a, 1990a, 1990b and 1992b, (90, 91, 92, 93 and Wenning and Ingersoll (2002 (94)) for a critique of these methods). The sediment assessment approaches listed in Table 2 can be classified as numeric (for example, equilibrium partitioning), descriptive (for example, whole-sediment toxicity tests), or a combination of numeric and descriptive approaches (for example, ERM, USEPA, 1992c, (95). Numeric methods can be used to derive chemical-specific sediment quality guidelines (SQGs). Descriptive methods such as toxicity tests with field-collected sediment cannot be used alone to develop numerical SQGs for individual chemicals. Although each approach can be used to make site-specific decisions, no one single approach can adequately address sediment quality. Overall, an integration of several methods using the weight of evidence is the most desirable approach for assessing the effects of contaminants associated with sediment, (Long et al. 1991(96) MacDonald et al. 1996 (97) Ingersoll et al. 1996 (98) Ingersoll et al. 1997 (99), Wenning and Ingersoll 2002 (94)). Hazard evaluations integrating data from laboratory exposures, chemical analyses, and benthic community assessments (the sediment quality triad) provide strong complementary evidence of the degree of pollution-induced degradation in aquatic communities (Burton, 1991 (69), Chapman 1992, 1997 (100, 101).)
Regulatory Applications—Test Method E 1706 provides information on the regulatory applications of sediment toxicity tests.
The USEPA Environmental Monitoring Management Council (EMMC) recommended the use of performance-based methods in developing standards, (Williams, 1993 (102). Performance-based methods were defined by EMMC as a monitoring approach which permits the use of appropriate methods that meet preestablished demonstrated performance standards (11.2).
The USEPA Office of Water, Office of Science and Technology, and Office of Research and Development held a workshop to provide an opportunity for experts in the field of sediment toxicology and staff from the USEPA Regional and Headquarters Program offices to discuss the development of standard freshwater, estuarine, and marine sediment testing procedures (USEPA, 1992a, 1994a (89, 103)). Workgroup participants arrived at a consensus on several culturing and testing methods. In developing guidance for culturing test organisms to be included in the USEPA methods manual for sediment tests, it was agreed that no one method should be required to culture organisms. However, the consensus at the workshop was that success of a test depends on the health of the cultures. Therefore, having healthy test organisms of known quality and age for testing was determined to be the key consideration relative to culturing methods. A performance-based criteria approach was selected in USEPA, 2000 (73) as the preferred method through which individual laboratories could use unique culturing methods rather than requiring use of one culturing method.
This standard recommends the use of performance-based criteria to allow each laboratory to optimize culture methods and minimize effects of test organism health on the reliability and comparability of test results. See Annex A1 and Annex A2 for a listing of performance criteria for culturing or testing.
1.1 This test method covers procedures for testing estuarine or marine organisms in the laboratory to evaluate the toxicity of contaminants associated with whole sediments. Sediments may be collected from the field or spiked with compounds in the laboratory. General guidance is presented in Sections 1-15 for conducting sediment toxicity tests with estuarine or marine amphipods. Specific guidance for conducting 10-d sediment toxicity tests with estuarine or marine amphipods is outlined in Annex A1 and specific guidance for conducting 28-d sediment toxicity tests with Leptocheirus plumulosus is outlined in Annex A2.
1.2 Procedures are described for testing estuarine or marine amphipod crustaceans in 10-d laboratory exposures to evaluate the toxicity of contaminants associated with whole sediments (Annex A1; USEPA 1994a (1)). Sediments may be collected from the field or spiked with compounds in the laboratory. A toxicity method is outlined for four species of estuarine or marine sediment-burrowing amphipods found within United States coastal waters. The species are Ampelisca abdita, a marine species that inhabits marine and mesohaline portions of the Atlantic coast, the Gulf of Mexico, and San Francisco Bay; Eohaustorius estuarius, a Pacific coast estuarine species; Leptocheirus plumulosus, an Atlantic coast estuarine species; and Rhepoxynius abronius, a Pacific coast marine species. Generally, the method described may be applied to all four species, although acclimation procedures and some test conditions (that is, temperature and salinity) will be species-specific (Sections 12 and Annex A1). The toxicity test is conducted in 1-L glass chambers containing 175 mL of sediment and 775 mL of overlying seawater. Exposure is static (that is, water is not renewed), and the animals are not fed over the 10-d exposure period. The endpoint in the toxicity test is survival with reburial of surviving amphipods as an additional measurement that can be used as an endpoint for some of the test species (for R. abronius and E. estuarius). Performance criteria established for this test include the average survival of amphipods in negative control treatment must be greater than or equal to 90 %. Procedures are described for use with sediments with pore-water salinity ranging from >0 o/ooto fully marine.
1.3 A procedure is also described for determining the chronic toxicity of contaminants associated with whole sediments with the amphipod Leptocheirus plumulosus in laboratory exposures (Annex A2; USEPA-USACE 2001(2)). The toxicity test is conducted for 28 d in 1-L glass chambers containing 175 mL of sediment and about 775 mL of overlying water. Test temperature is 25° ± 2°C, and the recommended overlying water salinity is 5 o/oo ± 2 o/oo(for test sediment with pore water at 1 o/oo to 10 o/oo) or 20 o/oo ± 2 o/oo (for test sediment with pore water >10 o/oo). Four hundred millilitres of overlying water is renewed three times per week, at which times test organisms are fed. The endpoints in the toxicity test are survival, growth, and reproduction of amphipods. Performance criteria established for this test include the average survival of amphipods in negative control treatment must be greater than or equal to 80 % and there must be measurable growth and reproduction in all replicates of the negative control treatment. This test is applicable for use with sediments from oligohaline to fully marine environments, with a silt content greater than 5 % and a clay content less than 85 %.
1.4 A salinity of 5 or 20 o/oo is recommended for routine application of 28-d test with L. plumulosus (Annex A2; USEPA-USACE 2001 (2)) and a salinity of 20 o/oois recommended for routine application of the 10-d test with E. estuarius or L. plumulosus (Annex A1). However, the salinity of the overlying water for tests with these two species can be adjusted to a specific salinity of interest (for example, salinity representative of site of interest or the objective of the study may be to evaluate the influence of salinity on the bioavailability of chemicals in sediment). More importantly, the salinity tested must be within the tolerance range of the test organisms (as outlined in Annex A1 and Annex A2). If tests are conducted with procedures different from those described in 1.3 or in Table A1.1 (for example, different salinity, lighting, temperature, feeding conditions), additional tests are required to determine comparability of results (1.10). If there is not a need to make comparisons among studies, then the test could be conducted just at a selected salinity for the sediment of interest.
1.5 Future revisions of this standard may include additional annexes describing whole-sediment toxicity tests with other groups of estuarine or marine invertebrates (for example, information presented in Guide E 1611 on sediment testing with polychaetes could be added as an annex to future revisions to this standard). Future editions to this standard may also include methods for conducting the toxicity tests in smaller chambers with less sediment (Ho et al. 2000 (3), Ferretti et al. 2002 (4)).
1.6 Procedures outlined in this standard are based primarily on procedures described in the USEPA (1994a (1)), USEPA-USACE (2001(2)), Test Method E 1706, and Guides E 1391, E 1525, E 1688, Environment Canada (1992 (5)), DeWitt et al. (1992a (6); 1997a (7)), Emery et al. (1997 (8)), and Emery and Moore (1996 (9)), Swartz et al. (1985 (10)), DeWitt et al. (1989 (11)), Scott and Redmond (1989 (12)), and Schlekat et al. (1992 (13)).
1.7 Additional sediment toxicity research and methods development are now in progress to (1) refine sediment spiking procedures, (2) refine sediment dilution procedures, (3) refine sediment Toxicity Identification Evaluation (TIE) procedures, (4) produce additional data on confirmation of responses in laboratory tests with natural populations of benthic organisms (that is, field validation studies), and (5) evaluate relative sensitivity of endpoints measured in 10- and 28-d toxicity tests using estuarine or marine amphipods. This information will be described in future editions of this standard.
1.8 Although standard procedures are described in Annex A2 of this standard for conducting chronic sediment tests with L. plumulosus, further investigation of certain issues could aid in the interpretation of test results. Some of these issues include further investigation to evaluate the relative toxicological sensitivity of the lethal and sublethal endpoints to a wide variety of chemicals spiked in sediment and to mixtures of chemicals in sediments from contamination gradients in the field (USEPA-USACE 2001 (2)). Additional research is needed to evaluate the ability of the lethal and sublethal endpoints to estimate the responses of populations and communities of benthic invertebrates to contaminated sediments. Research is also needed to link the toxicity test endpoints to a field-validated population model of L. plumulosus that would then generate estimates of population-level responses of the amphipod to test sediments and thereby provide additional ecologically relevant interpretive guidance for the laboratory toxicity test.
1.9 This standard outlines specific test methods for evaluating the toxicity of sediments with A. abdita, E. estuarius, L. plumulosus, and R. abronius. While standard procedures are described in this standard, further investigation of certain issues could aid in the interpretation of test results. Some of these issues include the effect of shipping on organism sensitivity, additional performance criteria for organism health, sensitivity of various populations of the same test species, and confirmation of responses in laboratory tests with natural benthos populations.
1.10 General procedures described in this standard might be useful for conducting tests with other estuarine or marine organisms (for example, Corophium spp., Grandidierella japonica, Lepidactylus dytiscus, Streblospio benedicti), although modifications may be necessary. Results of tests, even those with the same species, using procedures different from those described in the test method may not be comparable and using these different procedures may alter bioavailability. Comparison of results obtained using modified versions of these procedures might provide useful information concerning new concepts and procedures for conducting sediment tests with aquatic organisms. If tests are conducted with procedures different from those described in this test method, additional tests are required to determine comparability of results. General procedures described in this test method might be useful for conducting tests with other aquatic organisms; however, modifications may be necessary.
1.11 Selection of Toxicity Testing Organisms:
1.11.1 The choice of a test organism has a major influence on the relevance, success, and interpretation of a test. Furthermore, no one organism is best suited for all sediments. The following criteria were considered when selecting test organisms to be described in this standard (Table 1 and Guide E 1525). Ideally, a test organism should: (1) have a toxicological database demonstrating relative sensitivity to a range of contaminants of interest in sediment, (2) have a database for interlaboratory comparisons of procedures (for example, round-robin studies), (3) be in direct contact with sediment, (4) be readily available from culture or through field collection, (5) be easily maintained in the laboratory, (6) be easily identified, (7) be ecologically or economically important, (8) have a broad geographical distribution, be indigenous (either present or historical) to the site being evaluated, or have a niche similar to organisms of concern (for example, similar feeding guild or behavior to the indigenous organisms), (9) be tolerant of a broad range of sediment physico-chemical characteristics (for example, grain size), and (10) be compatible with selected exposure methods and endpoints (Guide E 1525). Methods utilizing selected organisms should also be (11) peer reviewed (for example, journal articles) and (12) confirmed with responses with natural populations of benthic organisms.
1.11.2 Of these criteria (Table 1), a database demonstrating relative sensitivity to contaminants, contact with sediment, ease of culture in the laboratory or availability for field-collection, ease of handling in the laboratory, tolerance to varying sediment physico-chemical characteristics, and confirmation with responses with natural benthic populations were the primary criteria used for selecting A. abdita, E. estuarius, L. plumulosus, and R. abronius for the current edition of this standard for 10-d sediment tests (Annex A1). The species chosen for this method are intimately associated with sediment, due to their tube- dwelling or free-burrowing, and sediment ingesting nature. Amphipods have been used extensively to test the toxicity of marine, estuarine, and freshwater sediments (Swartz et al., 1985 (10); DeWitt et al., 1989 (11); Scott and Redmond, 1989 (12); DeWitt et al., 1992a (6); Schlekat et al., 1992 (13)). The selection of test species for this standard followed the consensus of experts in the field of sediment toxicology who participated in a workshop entitled “Testing Issues for Freshwater and Marine Sediments”. The workshop was sponsored by USEPA Office of Water, Office of Science and Technology, and Office of Research and Development, and was held in Washington, D.C. from 16-18 September 1992 (USEPA, 1992 (14)). Of the candidate species discussed at the workshop, A. abdita, E. estuarius, L. plumulosus, and R. abronius best fulfilled the selection criteria, and presented the availability of a combination of one estuarine and one marine species each for both the Atlantic (the estuarine L. plumulosus and the marine A. abdita) and Pacific (the estuarine E. estuarius and the marine R. abronius) coasts. Ampelisca abdita is also native to portions of the Gulf of Mexico and San Francisco Bay. Many other organisms that might be appropriate for sediment testing do not now meet these selection criteria because little emphasis has been placed on developing standardized testing procedures for benthic organisms. For example, a fifth species, Grandidierella japonica was not selected because workshop participants felt that the use of this species was not sufficiently broad to warrant standardization of the method. Environment Canada (1992 (5)) has recommended the use of the following amphipod species for sediment toxicity testing: Amphiporeia virginiana, Corophium volutator, Eohaustorius washingtonianus, Foxiphalus xiximeus, and Leptocheirus pinguis. A database similar to those available for A. abdita, E. estuarius, L. plumulosus, and R. abronius must be developed in order for these and other organisms to be included in future editions of this standard.
1.11.3 The primary criterion used for selecting L. plumulosus for chronic testing of sediments was that this species is found in both oligohaline and mesohaline regions of estuaries on the East Coast of the United States and is tolerant to a wide range of sediment grain size distribution (USEPA-USACE 2001 (2), Annex Annex A2). This species is easily cultured in the laboratory and has a relatively short generation time (that is, about 24 d at 23°C, DeWitt et al. 1992a (6)) that makes this species adaptable to chronic testing (Section 12).
1.11.4 An important consideration in the selection of specific species for test method development is the existence of information concerning relative sensitivity of the organisms both to single chemicals and complex mixtures. Several studies have evaluated the sensitivities of A. abdita, E. estuarius, L. plumulosus, or R. abronius, either relative to one another, or to other commonly tested estuarine or marine species. For example, the sensitivity of marine amphipods was compared to other species that were used in generating saltwater Water Quality Criteria. Seven amphipod genera, including Ampelisca abdita and Rhepoxynius abronius, were among the test species used to generate saltwater Water Quality Criteria for 12 chemicals. Acute amphipod toxicity data from 4-d water-only tests for each of the 12 chemicals was compared to data for (1) all other species, (2) other benthic species, and (3) other infaunal species. Amphipods were generally of median sensitivity for each comparison. The average percentile rank of amphipods among all species tested was 57 %; among all benthic species, 56 %; and, among all infaunal species, 54 %. Thus, amphipods are not uniquely sensitive relative to all species, benthic species, or even infaunal species (USEPA 1994a (1)). Additional research may be warranted to develop tests using species that are consistently more sensitive than amphipods, thereby offering protection to less sensitive groups.
1.11.5 Williams et al. (1986 (15)) compared the sensitivity of the R. abronius 10-d whole sediment test, the oyster embryo (Crassostrea gigas) 48-h abnormality test, and the bacterium (Vibrio fisheri) 1-h luminescence inhibition test (that is, the Microtox test) to sediments collected from 46 contaminated sites in Commencement Bay, WA. Rhepoxynius abronius were exposed to whole sediment, while the oyster and bacterium tests were conducted with sediment elutriates and extracts, respectfully. Microtox was the most sensitive test, with 63 % of the sites eliciting significant inhibition of luminescence. Significant mortality of R. abronius was observed in 40 % of test sediments, and oyster abnormality occurred in 35 % of sediment elutriates. Complete concordance (that is, sediments that were either toxic or not-toxic in all three tests) was observed in 41 % of the sediments. Possible sources for the lack of concordance at other sites include interspecific differences in sensitivity among test organisms, heterogeneity in contaminant types associated with test sediments, and differences in routes of exposure inherent in each toxicity test. These results highlight the importance of using multiple assays when performing sediment assessments.
1.11.6 Several studies have compared the sensitivity of combinations of the four amphipods to sediment contaminants. For example, there are several comparisons between A. abdita and R. abronius, between E. estuarius and R. abronius, and between A. abdita and L. plumulosus. There are fewer examples of direct comparisons between E. estuarius and L. plumulosus, and no examples comparing L. plumulosus and R. abronius. There is some overlap in relative sensitivity from comparison to comparison within each species combination, which appears to indicate that all four species are within the same range of relative sensitivity to contaminated sediments.
220.127.116.11 Word et al. (1989 (16)) compared the sensitivity of A. abdita and R. abronius to contaminated sediments in a series of experiments. Both species were tested at 15°C. Experiments were designed to compare the response of the organism rather than to provide a comparison of the sensitivity of the methods (that is, Ampelisca abdita would normally be tested at 20°C). Sediments collected from Oakland Harbor, CA, were used for the comparisons. Twenty-six sediments were tested in one comparison, while 5 were tested in the other. Analysis of results using Kruskal Wallace rank sum test for both experiments demonstrated that R. abronius exhibited greater sensitivity to the sediments than A. abdita at 15°C. Long and Buchman (1989 (17)) also compared the sensitivity of A. abdita and R. abronius to sediments from Oakland Harbor, CA. They also determined that A. abdita showed less sensitivity than R. abronius, but they also showed that A. abdita was less sensitive to sediment grain size factors than R. abronius.
18.104.22.168 DeWitt et al. (1989 (11)) compared the sensitivity of E. estuarius and R. abronius to sediment spiked with fluoranthene and field-collected sediment from industrial waterways in Puget Sound, WA, in 10-d tests, and to aqueous cadmium (CdCl2) in a 4-d water-only test. The sensitivity of E. estuarius was from two (to spiked-spiked sediment) to seven (to one Puget Sound, WA, sediment) times less sensitive than R. abronius in sediment tests, and ten times less sensitive to CdCl2 in the water-only test. These results are supported by the findings of Pastorok and Becker (1990 (18)) who found the acute sensitivity of E. estuarius and R. abronius to be generally comparable to each other, and both were more sensitive than Neanthes arenaceodentata (survival and biomass endpoints), Panope generosa (survival), and Dendraster excentricus (survival).
22.214.171.124 Leptocheirus plumulosus was as sensitive as the freshwater amphipod Hyalella azteca to an artificially created gradient of sediment contamination when the latter was acclimated to oligohaline salinity (that is, 6 o/oo; McGee et al., 1993 (19)). DeWitt et al. (1992b (20)) compared the sensitivity of L. plumulosus with three other amphipod species, two mollusks, and one polychaete to highly contaminated sediment collected from Baltimore Harbor, MD, that was serially diluted with clean sediment. Leptocheirus plumulosus was more sensitive than the amphipods Hyalella azteca and Lepidactylus dytiscus and exhibited equal sensitivity with E. estuarius. Schlekat et al. (1995 (21)) describe the results of an interlaboratory comparison of 10-d tests with A. abdita, L. plumulosus and E. estuarius using dilutions of sediments collected from Black Rock Harbor, CT. There was strong agreement among species and laboratories in the ranking of sediment toxicity and the ability to discriminate between toxic and non-toxic sediments.
126.96.36.199 Hartwell et al. (2000 (22)) evaluated the response of Leptocheirus plumulosus (10-d survival or growth) to the response of the amphipod Lepidactylus dytiscus (10-d survival or growth), the polychaete Streblospio benedicti (10-d survival or growth), and lettuce germination (Lactuca sativa in 3-d exposure) and observed that L. plumulosus was relatively insensitive compared to the response of either L. dytiscus or S. benedicti in exposures to 4 sediments with elevated metal concentrations.
188.8.131.52 Ammonia is a naturally occurring compound in marine sediment that results from the degradation of organic debris. Interstitial ammonia concentrations in test sediment can range from <1 mg/L to in excess of 400 mg/L (Word et al., 1997 (23)). Some benthic infauna show toxicity to ammonia at concentrations of about 20 mg/L (Kohn et al., 1994 (24)). Based on water-only and spiked-sediment experiments with ammonia, threshold limits for test initiation and termination have been established for the L. plumulosus chronic test. Smaller (younger) individuals are more sensitive to ammonia than larger (older) individuals (DeWitt et al., 1997a (7), b (25). Results of a 28-d test indicated that neonates can tolerate very high levels of pore-water ammonia (>300 mg/L total ammonia) for short periods of time with no apparent long-term effects (Moore et al., 1997 (26)). It is not surprising L. plumulosus has a high tolerance for ammonia given that these amphipods are often found in organic rich sediments in which diagenesis can result in elevated pore-water ammonia concentrations. Insensitivity to ammonia by L. plumulosus should not be construed as an indicator of the sensitivity of the L. plumulosus sediment toxicity test to other chemicals of concern.
1.11.7 Limited comparative data is available for concurrent water-only exposures of all four species in single-chemical tests. Studies that do exist generally show that no one species is consistently the most sensitive.
184.108.40.206 The relative sensitivity of the four amphipod species to ammonia was determined in ten-d water only toxicity tests in order to aid interpretation of results of tests on sediments where this toxicant is present (USEPA 1994a (1)). These tests were static exposures that were generally conducted under conditions (for example, salinity, photoperiod) similar to those used for standard 10-d sediment tests. Departures from standard conditions included the absence of sediment and a test temperature of 20°C for L. plumulosus, rather than 25°C as dictated in this standard. Sensitivity to total ammonia increased with increasing pH for all four species. The rank sensitivity was R. abronius = A. abdita > E. estuarius > L. plumulosus. A similar study by Kohn et al. (1994 (24)) showed a similar but slightly different relative sensitivity to ammonia with A. abdita > R. abronius = L. plumulosus > E. estuarius.
220.127.116.11 Cadmium chloride has been a common reference toxicant for all four species in 4-d exposures. DeWitt et al. (1992a (6)) reports the rank sensitivity as R. abronius > A. abdita > L. plumulosus > E. estuarius at a common temperature and salinity of 15°C and 28 o/oo. A series of 4-d exposures to cadmium that were conducted at species-specific temperatures and salinities showed the following rank sensitivity: A. abdita = L. plumulosus = R. abronius > E. estuarius (USEPA 1994a (1)).
18.104.22.168 Relative species sensitivity frequently varies among contaminants; consequently, a battery of tests including organisms representing different trophic levels may be needed to assess sediment quality (Craig, 1984 (27); Williams et al. 1986 (15); Long et al., 1990 (28); Ingersoll et al., 1990 (29); Burton and Ingersoll, 1994 (31)). For example, Reish (1988 (32)) reported the relative toxicity of six metals (arsenic, cadmium, chromium, copper, mercury, and zinc) to crustaceans, polychaetes, pelecypods, and fishes and concluded that no one species or group of test organisms was the most sensitive to all of the metals.
1.11.8 The sensitivity of an organism is related to route of exposure and biochemical response to contaminants. Sediment-dwelling organisms can receive exposure from three primary sources: interstitial water, sediment particles, and overlying water. Food type, feeding rate, assimilation efficiency, and clearance rate will control the dose of contaminants from sediment. Benthic invertebrates often selectively consume different particle sizes (Harkey et al. 1994 (33)) or particles with higher organic carbon concentrations which may have higher contaminant concentrations. Grazers and other collector-gatherers that feed on aufwuchs and detritus may receive most of their body burden directly from materials attached to sediment or from actual sediment ingestion. In some amphipods (Landrum, 1989 (34)) and clams (Boese et al., 1990 (35)) uptake through the gut can exceed uptake across the gills for certain hydrophobic compounds. Organisms in direct contact with sediment may also accumulate contaminants by direct adsorption to the body wall or by absorption through the integument (Knezovich et al. 1987 (36)).
1.11.9 Despite the potential complexities in estimating the dose that an animal receives from sediment, the toxicity and bioaccumulation of many contaminants in sediment such as Kepone®, fluoranthene, organochlorines, and metals have been correlated with either the concentration of these chemicals in interstitial water or in the case of non-ionic organic chemicals, concentrations in sediment on an organic carbon normalized basis (Di Toro et al. 1990 (37); Di Toro et al. 1991(38)). The relative importance of whole sediment and interstitial water routes of exposure depends on the test organism and the specific contaminant (Knezovich et al. 1987 (36)). Because benthic communities contain a diversity of organisms, many combinations of exposure routes may be important. Therefore, behavior and feeding habits of a test organism can influence its ability to accumulate contaminants from sediment and should be considered when selecting test organisms for sediment testing.
1.11.10 The use of A. abdita, E. estuarius, R. abronius, and L. plumulosus in laboratory toxicity studies has been field validated with natural populations of benthic organisms (Swartz et al. 1994 (39) and Anderson et al. 2001 (40) for E. estuarius, Swartz et al. 1982 (43) and Anderson et al. 2001 (40) for R. abronius, McGee et al. 1999 (41)and McGee and Fisher 1999 (42) for L. plumulosus).
22.214.171.124 Data from USEPA Office of Research and Development's Environmental Monitoring and Assessment program were examined to evaluate the relationship between survival of Ampelisca abdita in sediment toxicity tests and the presence of amphipods, particularly ampeliscids, in field samples. Over 200 sediment samples from two years of sampling in the Virginian Province (Cape Cod, MA, to Cape Henry, VA) were available for comparing synchronous measurements of A. abdita survival in toxicity tests to benthic community enumeration. Although species of this genus were among the more frequently occurring taxa in these samples, ampeliscids were totally absent from stations that exhibited A. abdita test survival <60 % of that in control samples. Additionally, ampeliscids were found in very low densities at stations with amphipod test survival between 60 and 80 % (USEPA 1994a (1)). These data indicate that tests with
2. Referenced Documents (purchase separately) The documents listed below are referenced within the subject standard but are not provided as part of the standard.
D1129 Terminology Relating to Water
D4447 Guide for Disposal of Laboratory Chemicals and Samples
E29 Practice for Using Significant Digits in Test Data to Determine Conformance with Specifications
E105 Practice for Probability Sampling of Materials
E122 Practice for Calculating Sample Size to Estimate, With Specified Precision, the Average for a Characteristic of a Lot or Process
E141 Practice for Acceptance of Evidence Based on the Results of Probability Sampling
E177 Practice for Use of the Terms Precision and Bias in ASTM Test Methods
E178 Practice for Dealing With Outlying Observations
E456 Terminology Relating to Quality and Statistics
E691 Practice for Conducting an Interlaboratory Study to Determine the Precision of a Test Method
E729 Guide for Conducting Acute Toxicity Tests on Test Materials with Fishes, Macroinvertebrates, and Amphibians
E943 Terminology Relating to Biological Effects and Environmental Fate
E1241 Guide for Conducting Early Life-Stage Toxicity Tests with Fishes
E1325 Terminology Relating to Design of Experiments
E1391 Guide for Collection, Storage, Characterization, and Manipulation of Sediments for Toxicological Testing and for Selection of Samplers Used to Collect Benthic Invertebrates
E1402 Guide for Sampling Design
E1525 Guide for Designing Biological Tests with Sediments
E1611 Guide for Conducting Sediment Toxicity Tests with Polychaetous Annelids
E1688 Guide for Determination of the Bioaccumulation of Sediment-Associated Contaminants by Benthic Invertebrates
E1706 Test Method for Measuring the Toxicity of Sediment-Associated Contaminants with Freshwater Invertebrates
E1847 Practice for Statistical Analysis of Toxicity Tests Conducted Under ASTM Guidelines
E1850 Guide for Selection of Resident Species as Test Organisms for Aquatic and Sediment Toxicity Tests
Ampelisca abdita; amphipod; bioavailability; chronic; Eohaustorius estuarius; estuarine; invertebrates; Leptocheirus plumulosus; marine; Rhepoxynius abronius; sediment; toxicity; Acidity, alkalinity, pH--chemicals; Acute toxicity tests; Ampelisca abdita; Amphipods/Amphibia; Aqueous environments; Benthic macroinvertebrates (collecting); Biological data analysis--sediments; Bivalve molluscs; Chemical analysis--water applications; Contamination--environmental; Corophium; Crustacea; EC50 test; Eohaustorius estuarius; Estuarine environments; Field testing--environmental materials/applications; Geochemical characteristics; Grandidierella japonica; Leptocheirus Plumuulosus; Marine environments; Median lethal dose; Polychaetes; Reference toxicants; Rhepoxynium abronius; Saltwater; Seawater (natural/synthetic); Sediment toxicity testing; Static tests--environmental materials/applications; Ten-day testing; Toxicity/toxicology--water environments
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Water reuse can be defined as the use of reclaimed water for a direct beneficial purpose. The use of reclaimed water for irrigation and other purposes has been employed as a water conservation practice in Florida, California, Texas, Arizona, and other states for many years.
Reclaimed water, also known as recycled water, is water recovered from domestic, municipal, and industrial wastewater treatment plants that has been treated to standards that allow safe reuse. Properly reclaimed water is typically safe for most uses except human consumption.
Wastewater is not reclaimed water. Wastewater is untreated liquid industrial waste and/or domestic sewage from residential dwellings, commercial buildings, and industrial facilities. Gray water, or untreated wastewater from bathing or washing, is one form of wastewater. Wastewater may be land applied, but this is considered to be land treatment rather than water reuse.
The demand for fresh water in Virginia is growing as the state’s population increases. This demand can potentially exceed supply during times of even moderate drought. In recent years, the normal seasonal droughts that have occurred in Virginia have caused local and state government to enact water conservation ordinances. These ordinances limit the use of potable water (water suitable for human consumption) for such things as car washing and landscape irrigation. The potential for developing new sources of potable water is limited. Conservation measures, such as irrigating with reclaimed water, are one way to help ensure existing water supplies are utilized as efficiently as possible.
The environmental benefits of using reclaimed water include:
Reclaimed water typically comes from municipal wastewater treatment plants, although some industries (e.g., food processors) also generate water that may be suitable for nonpotable uses. (Figure 1).
During primary treatment at a wastewater treatment plant, inorganic and organic suspended solids are removed from plant influent by screening, and settling. The decanted effluent from the primary treatment process is then subjected to secondary treatment, which involves biological decomposition of organic material and settling to further separate water from solids. If a wastewater treatment plant is not equipped to perform advanced treatment, water is disinfected and discharged to natural water bodies following secondary treatment.
Advanced or tertiary treatment consists of further removal of suspended and dissolved solids, including nutrients, and disinfection. Advanced treatment can include:
Water that has undergone advanced treatment is disinfected prior to being released or reused. Reclaimed water often requires greater treatment than effluent that is discharged to local streams or rivers because users will typically have more direct contact with undiluted reclaimed water than undiluted effluent.
For an interactive diagram of a wastewater treatment system with more information on treatment processes, please see www.wef.org/apps/gowithflow/theflow.htm.
Although the primary focus of this publication is on the use of reclaimed water for agricultural, municipal, and residential irrigation, reclaimed water can be used for many other purposes. Non-irrigation uses for reclaimed water include:
Intentional indirect potable reuse means that reclaimed water is discharged to a water body where it is then purposefully used as a raw water supply for another water treatment plant. This occurs unintentionally in most rivers, since downstream water treatment plants use treated water discharged by upstream wastewater treatment plants.
Direct potable reuse refers to the use of reclaimed water for drinking directly after treatment, and, to date, has only been implemented in Africa (U.S. EPA, 2004).
Examples of non-irrigation permitted water reuse projects in Virginia are:
The turfgrass and ornamental horticulture industries have grown as Virginia becomes more urbanized. The acreage devoted to high-value specialty crops that benefit from irrigation, such as fruits and vegetables, is also increasing. As demand for potable water increases, maintaining turf, landscape plants, and crops will require the utilization of previously underutilized water sources.
The regulation of reclaimed water production and use encourages both the supply of and the demand for reclaimed water. The benefits to suppliers of reclaimed water include greater public awareness and demand for reclaimed water and clear guidelines for reclaimed water production. Benefits to end users include increased public acceptance of the use of reclaimed water and a subsequent decrease in the demand for fresh water.
There are no federal regulations governing reclaimed water use, but the U.S. EPA (2004) has established guidelines to encourage states to develop their own regulations. The primary purpose of federal guidelines and state regulations is to protect human health and water quality. To reduce disease risks to acceptable levels, reclaimed water must meet certain disinfection standards by either reducing the concentrations of constituents that may affect public health and/or limiting human contact with reclaimed water.
The U.S. EPA (2004) recommends that water intended for reuse should:
Biochemical oxygen demand (BOD) is an indicator of the presence of reactive organic matter in water. Total suspended solids (TSS) or turbidity (measured in nephelometric turbidity units, or NTUs) are measures of the amount of organic and inorganic particulate matter in water. Some other parameters often measured as indicators of disinfection efficiency include:
The recommended values for each of these indicators depend on the intended use of the reclaimed water (Table 1).
Table 1. Summary of U.S. EPA guidelines for water reuse for irrigation
(Adapted from U.S. EPA, 2004).
Monitoring for specific pathogens and microconstituents may become a part of the standard testing protocol as the use of reclaimed water for indirect potable reuse applications increases. Pathogens of particular concern include enteric viruses and the protozoan parasites Giardia and Cryptosporidium, whose monitoring is required by the state of Florida for water reuse projects.
Microconstituents include organic chemicals, such as pharmaceutically active substances, personal care products, endocrine disrupting compounds, and previously unregulated inorganic elements whose toxicity may be re-assessed or newly evaluated. Fish, amphibians, and birds have been found to develop reproductive system abnormalities upon direct or indirect exposure to a variety of endocrine disrupting compounds. Such microconstituents may have the potential to cause reproduction system abnormalities and immune system malfunctioning in other wildlife and humans at higher concentrations. The impacts of the extremely low concentrations of these compounds found in wastewater effluent or reclaimed water are unknown. To date, there is no evidence that microconstituents cause human health effects at environmentally relevant concentrations.
Some possible options for the removal of microconstituents from wastewater are treatment with ozone, hydrogen peroxide, and UV light. These methods can destroy some microconstituents via advanced oxidation, but the endocrine disruption activity of the by-products created during oxidation may also be of concern.
No illnesses have been directly associated with the use of properly treated reclaimed water in the U.S. (U.S. EPA, 2004). The U.S. EPA recommends, however, that ongoing research and additional monitoring for Giardia, Cryptosporidium, and microconstituents be conducted to understand changes in reclaimed water quality.
State regulations need not agree with U.S. EPA guidelines and are often more stringent. In Virginia, water reuse means direct beneficial reuse, indirect potable reuse, or a controlled use in accordance with the Water Reclamation and Reuse Regulation (9 VAC 25-740-10 et seq.; available at the Virginia Department of Environmental Quality website www.deq.virginia.gov/programs/homepage.html under Water Reuse and Reclamation.)
The Virginia Water Regulation and Reuse Regulation establishes legal requirements for the reclamation and treatment of water that is to be reused. These require ments are designed to protect both water quality and public health, while encouraging the use of reclaimed water. The Virginia Department of Environmental Quality, Water Quality Division has oversight over the Virginia Water Reclamation and Reuse Regulation.
The primary determinants of how reclaimed water of varying quality can be used are based on treatment processes to which the water has been subjected and on quantitative chemical, physical, and biological standards. Reclaimed water suitable for reuse in Virginia is categorized as either Level 1 or Level 2 (Table 2). The minimum standard requirements for reclaimed water for specific uses are summarized in Table 3.
Table 2. Minimum standards for treatment of Level 1 and Level 2 reclaimed water.
(Summarized from Virginia Water Reclamation and Reuse Regulations: 9 VAC 25-740-10 et seq.)
Table 3. Minimum treatment requirements for irrigation and landscape-related reuse of reclaimed water in Virginia.
(Summarized from Virginia Water Reclamation and Reuse Regulations: 9 VAC 25-740-10 et seq.)
Water quality must be considered when using reclaimed water for irrigation. The following properties are critical to plant and soil health and environmental quality.
Salinity, or salt concentration, is probably the most important consideration in determining whether water is suitable for reuse (U.S. EPA, 2004). Water salinity is the sum of all elemental ions (e.g., sodium, calcium, chloride, boron, sulfate, nitrate) and is usually measured by determining the electrical conductivity (EC, units = dS/m) or total dissolved solids (TDS, units = mg/L) concentration of the water. Water with a TDS concentration of 640 mg/L will typically have an EC of approximately 1 dS/m.
Salts in reclaimed water come from:
Most reclaimed water from urban areas is slightly saline (TDS ≤ 1280 mg/L or EC ≤ 2 dS/m). High salt concentrations reduce water uptake in plants by lowering the osmotic potential of the soil. For instance, residential use of water adds approx 200-400 mg/L dissolved salts (Lazarova et al., 2004a). Plants differ in their sensitivity to salt levels so the salinity of the particular reclaimed water source should be measured so that appropriate crops and/or application rates can be selected. Most turfgrasses can tolerate water with 200-800 mg/L soluble salts, but salt levels above 2,000 mg/L may be toxic (Harivandi, 2004). For further information on managing turfgrasses when irrigating with saline water, see Carrow and Duncan (1998).
Many other crop and landscape plants are more sensitive to high soluble-salt levels than turfgrasses, and should be managed accordingly. See Wu and Dodge (2005) for a list of landscape plants with their relative salt tolerance and Maas (1987) for information on salt-tolerant crops.
Specific dissolved ions may also affect irrigation water quality. For example, irrigation water with a high concentration of sodium (Na) ions may cause dispersion of soil aggregates and sealing of soil pores. This is a particular problem in golf course irrigation (Sheikh, 2004) since soil compaction is already a concern due to persistent foot and vehicular traffic. The Sodium Adsorption Ratio (SAR), which measures the ratio of sodium to other ions, is used to evaluate the potential effect of irrigation water on soil structure. For more information on how to assess and interpret SAR levels, please see Harivandi (1999).
High levels of sodium can also be directly toxic to plants both through root uptake and by accumulation in plant leaves following sprinkler irrigation. The specific concentration of sodium that is considered to be toxic will vary with plant species and the type of irrigation system. Turfgrasses are generally more tolerant to sodium than most ornamental plant species.
Although boron (B) and chlorine (Cl) are necessary at low levels for plant growth, dissolved boron and chloride ions can cause toxicity problems at high concentrations. Specific toxic concentrations will vary depending on plant species and type of irrigation method used. Levels of boron as low as 1 to 2 mg/L in irrigation water can cause leaf burn on ornamental plants, but turfgrasses can often tolerate levels as high as 10 mg/L (Harivandi, 1999). Very salt-sensitive landscape plants such as crape myrtle (Lagerstroemia sp.), azalea (Rhododendron sp.), and Chinese privet (Ligustrum sinense) may be damaged by overhead irrigation with reclaimed water containing chloride levels over 100 mg/L, but most turfgrasses are relatively tolerant to chloride if they are mowed frequently (Harivandi, 1999; Crook, 2005).
Reclaimed water typically contains more nitrogen (N) and phosphorus (P) than drinking water. The amounts of N and P provided by the reclaimed water can be calculated as the product of the estimated irrigation volume and the N and P concentration in the water. To prevent N and P leaching into groundwater, the Virginia Water Reclamation and Reuse Regulation requires that a nutrient management plan be written for bulk use of reclaimed water not treated to achieve biological nutrient removal (BNR), which the regulation defines as treatment that achieves an annual average of 8.0 mg/L total N and 1.0 mg/L total P. Water that has been subjected to BNR treatment processes contains such low concentrations of N and P that the reclaimed water can be applied at rates sufficient to supply a crop’s water needs without risk of surface or ground water contamination.
The Virginia Water Reclamation and Reuse Regulations require that irrigation with reclaimed water shall be limited to supplemental irrigation. Supplemental irrigation is defined as that amount of water which, in combination with rainfall, meets the water demands of the irrigated vegetation to maximize production or optimize growth.
Irrigation rates for reclaimed water are site- and crop-specific, and will depend on the following factors (U.S. EPA, 2004; Lazarova et al., 2004b).
1. First, seasonal irrigation demands must be determined. These can be predicted with:• an evapotranspiration estimate for the particular crop being grown
• determination of the period of plant growth
• average annual precipitation data
• data for soil permeability and water holding capacity
Methods for calculating such irrigation requirements can be found in the U.S. Department of Agriculture’s National Engineering Handbook at www.info.usda.gov/CED/ftp/CED/neh-15.htm (USDA-NRCS, 2003) and in Reed et al. (1995). These calculations are more complicated for landscape plantings than for agricultural crops or turf because landscape plantings consist of many different species with different requirements.
2. The properties of the specific reclaimed water to be used, as detailed in the section above, must be taken into account since these may limit the total amount of water that can be applied per season.
3. The availability of the reclaimed water should also be quantified, including:• the total amount available
• the time of year when available
• availability of water storage facilities for the nongrowing season
• delivery rate and type
Water reuse is actively promoted by the Florida Department of Environmental Protection since Florida law requires that the use of potable water for irrigation be limited. In 2005, 462 Florida golf courses, covering over 56,000 acres of land, were irrigated with reclaimed water. Reclaimed water was also used to irrigate 201,465 residences, 572 parks, and 251 schools. St. Petersburg is home to one of the largest dual distribution systems in the world. (A dual distribution system is one where pipes carrying reclaimed water are separate from those carrying potable water.) In existence since the 1970s, this network provides reclaimed water to residences, golf courses, parks, schools, and commercial areas for landscape irrigation, and to commercial and industrial customers for cooling and other applications.
For more information, see Crook (2005) and Florida Department of Environmental Protection (2006).
The town of Cary is the first city in the state of North Carolina to institute a dual distribution system. The system has been in operation since 2001 and can provide up to 1 million gallons of reclaimed water daily for irrigating and cooling. The reclaimed water has undergone advanced treatment and meets North Carolina water quality rules. To date, there are over 400 residential and industrial users.
For more information, see www.townofcary.org/depts/pwdept/reclaimhome.htm.
The Bayberry Hills Golf Course expansion is one of numerous water reuse projects in Massachusetts. It was initiated in 2001 as an addition to an existing golf course of seven holes irrigated with reclaimed water. These seven holes use approximately 18 million gallons of water per year, and water reuse was necessary since Yarmouth’s water supply was already operating at capacity during summer months. The reused water undergoes secondary treatment followed by ozone treatment, filtration, and UV disinfection. There are provisions for water storage during the nongrowing season. The water reuse project has reduced the nitrogen needed for golf course fertilization.
For more information on this and other reuse projects in the state of Massachusetts, see www.mapc.org/regional_planning/MAPC_Water_Reuse_Report_2005.pdf. For further information on irrigation of golf courses with reclaimed water, see United States Golf Association (1994).
The Southeast Farm in Tallahassee, Florida, has been irrigating with reclaimed water since 1966. The farm is a cooperative between the city of Tallahassee, which supplies water, and farmers who contract acreage. Until 1980, the farm was limited to 20 acres of land for hay production, but has expanded since then to 2,163 acres. The irrigation water receives secondary treatment. The crops grown are corn (Zea mays L. subsp. Mays), soybeans [Glycine max (L.) Merr], bermudagrass [Cynodon dactylon (L.) Pers], and rye (Secale cereale L.).
In recent years, however, elevated nitrate levels have been found in the waters of Wakulla Springs State Park south of Tallahassee, which is one of the largest and deepest freshwater springs in the world. This has apparently resulted in excessive growth of algae and exotic aquatic plant species, causing reduced clarity and changes in the spring’s ecosystem. Dye studies have confirmed that at least a portion of the nitrate comes from the Southeast Farm’s irrigated fields, although studies are on-going. As a result, in June 2006, the city of Tallahassee removed all cattle from Southeast Farm, eliminated regular use of nitrogen fertilizer on the farm, and implemented a comprehensive nutrient management plan for the farm.
For more information, see www.talgov.com/you/water/pdf/sefarm.pdf or U.S. EPA (2004).
Water Conserv II has been in existence since 1986, and is the first project permitted by the Florida Department of Environmental Protection for crops for human consumption. Over 3,000 acres of citrus groves are irrigated with reclaimed water, in addition to nurseries, residential landscaping, a sand mine, and the Orange County National Golf Center. No problems have resulted from the irrigation. The reclaimed water provides adequate boron and phosphorus and maintains soil at correct pH for citrus growth. The adequate supply of water permits citrus growers to maintain optimum moisture levels for high yields and ample water for freeze protection, which requires more than eight times as much water as normal irrigation.
Although Water Conserv II had historically provided reclaimed water to citrus growers for no charge, the project recently began charging for water. It’s unclear if citrus growers will continue to irrigate with reclaimed water, or whether Water Conserv II’s emphasis will change to providing reclaimed water for residential, industrial, and landscape customers.
For more information, see www.waterconservii.com/ or U.S. EPA (2004).
This publication was reviewed by Adria Bordas, Bobby Clark, Erik Ervin, and Gary Felton. A draft version was reviewed by Bob Angelotti, Marcia Degen, Karen Harr, George Kennedy, Valerie Rourke, and Terry Wagner. Any opinions, conclusions, or recommendations expressed in this publication are those of the authors.
www.watereuse.org/: WateReuse Association. “The WateReuse Association is a non-profit organization whose mission is to advance the beneficial and efficient use of water resources through education, sound science, and technology using reclamation, recycling, reuse, and desalination for the benefit of our members, the public, and the environment.” Page contains links to water reuse projects (mostly in the western U.S.), and other useful links.
www.cvco.org/science/vwea/navbuttons/Glossary-11-01.pdf: Virginia Water Environment Association’s Virginia Water Reuse Glossary.
www.hrsd.com/waterreuse.htm: Hampton Roads (Virginia) Sanitation District water reuse page. Description of industrial water reuse project, research reports, FAQ’s, and glossary of water reuse jargon.
www.floridadep.org/water/reuse/index.htm: Florida Department of Environmental Protection water re-use page. Links to many water reuse-related resources on site, including general education/information materials, and Florida-specific links on water reuse policy, regulations, and projects.
www.gaepd.org/Files_PDF/techguide/wpb/reuse.pdf: Georgia Department of Natural Resources Environmental Protection Division’s “Guidelines for Water Reclamation and Urban Water Re-Use (2002).
www.mass.gov/dep/water/wastewater/wrfaqs.htm: Massachusetts Department of Environmental Protection FAQ on water reuse.
www.bcua.org/WPC_VT_WasteWaterReUse.htm: Bergen County (New Jersey) Utilities Authority. Describes reuse of wastewater effluent re-use in cooling towers and for sewer cleaning.
www.owasa.org/pages/WaterReuse/questionsandanswers.html: FAQ about Orange Water and SewerAuthority’s (Carrboro, NC) water reuse project for the University of North Carolina at Chapel Hill.
Carrow, R.N. and R.R. Duncan. 1998. Salt-affected turfgrass sites: Assessment and management. John Wiley & Sons, Inc., New York, N.Y.
Crook, James. 2005. St. Petersburg, Florida, dual water system: A case study. Water conservation, reuse, and recycling: Proceedings of an Iranian-American workshop. The National Academies Press, Washington, D.C.
Florida Department of Environmental Protection. 2006. 2005 reuse inventory. FDEP, Tallahassee, FL. Available on-line at www.floridadep.org/water/reuse/inventory.htm.
Harivandi, M. Ali. 1999. Interpreting turfgrass irrigation water test results. Publication 8009. University of California Division of Agriculture and Natural Resources, Oakland, Calif. Available on-line at anrcatalog.ucdavis.edu/pdf/8009.pdf.
Harivandi, M. Ali. 2004. Evaluating recycled waters for golf course irrigation. U.S. Golf Association Green Section Record 42(6): 25-29. Available on-line at turf.lib.msu.edu/2000s/2004/041125.pdf.
Landschoot, Peter. 2007. Irrigation water quality guidelines for turfgrass sites. Department of Crop and Soil Sciences, Cooperative Extension. Penn State University, State College, Pa. Available on-line at turfgrassmanagement.psu.edu/irrigation_water_quality_for_turfgrass_sites.cfm.
Lazarova, Valentina and Takashi Asano. 2004. Challenges of sustainable irrigation with recycled water. p. 1-30. In Valentina Lazarova and Bahri (ed.). Water reuse for irrigation: agriculture, landscapes, and turf grass. CRC Press, Boca Raton, Fla.
Lazarova, Valentina, Herman Bouwer, and Akica Bahri. 2004a. Water quality considerations. p. 31-60. In Valentina Lazarova and Akica Bahri (ed.). Water reuse for irrigation: agriculture, landscapes, and turf grass. CRC Press, Boca Raton, Fla.
Lazarova, Valentina, Ioannis Papadopoulous, and Akica Bahri. 2004b. Code of successful agronomic practices. p. 103-150. In Valentina Lazarova and Akica Bahri (ed.). Water reuse for irrigation: agriculture, landscapes, and turf grass. CRC Press, Boca Raton, Fla.
Maas, E.V. 1987. Salt tolerance of plants. p. 57–75. In B.R. Christie (ed.) CRC handbook of plant science in agriculture, Vol. II. CRC Press, Boca Raton, Fla.
Metropolitan Area Planning Council. 2005. Once is not enough: A guide to water reuse in Massachusetts. MAPC, Boston, Mass. Available on-line at www.mapc.org/regional_planning/MAPC_Water_Reuse_Report_2005.pdf.
Reed, Sherwood C., Ronald W. Crites, and E. Joe Middlebrooks. 1995. Natural systems for waste management and treatment. 2nd edition. McGraw-Hill, Inc. New York, N.Y.
Sheikh, Bahman. 2004. Code of practices for landscape and golf course irrigation. In Valentina Lazarova and Akica Bahri (ed.). Water reuse for irrigation: agriculture, landscapes, and turf grass. CRC Press, Boca Raton, Fla.
USDA-NRCS. 2003. Irrigation water requirements. Section 15, Chapter 2. p. 2-i-2-284. In Part 623 National Engineering Handbook. U.S. Dept. of Agriculture Natural Resources Conservation Service, Washington, D.C. Available on-line at www.info.usda.gov/CED/ftp/CED/neh-15.htm.
U.S. EPA. 2003. National primary drinking water standards. EPA 816-F-03-016. U.S. Environmental Protection Agency, Washington, D.C.
U.S. EPA. 2004. Guidelines for water reuse. EPA 645-R-04-108. U.S. Environmental Protection Agency, Washington, D.C. Available on-line at www.epa.gov/ORD/NRMRL/pubs/625r04108/625r04108.pdf.
United States Golf Association. 1994. Wastewater reuse for golf course irrigation. Lewis Publishers, Chelsea, Mich. 294 p.
VAAWW-VWEA. 2000. A Virginia water reuse glossary. Virginia Section, American Water Works Association and Virginia Water Environment Federation. Available on-line at www.cvco.org/science/vwea/navbuttons/Glossary-11-01.pdf.
Wu, Lin, and Linda Dodge. 2005. Landscape plant salt tolerance guide for recycled water irrigation. Slosson Research Endowment for Ornamental Horticulture, Department of Plant Sciences, University of California, Davis, Calif. Available on-line at ucce.ucdavis.edu/files/filelibrary/5505/20091.pdf.
Reviewed by Greg Evanylo, Extension Specialist, Crop and Soil Environmental Sciences
Virginia Cooperative Extension materials are available for public use, re-print, or citation without further permission, provided the use includes credit to the author and to Virginia Cooperative Extension, Virginia Tech, and Virginia State University.
Issued in furtherance of Cooperative Extension work, Virginia Polytechnic Institute and State University, Virginia State University, and the U.S. Department of Agriculture cooperating. Alan L. Grant, Dean, College of Agriculture and Life Sciences; Edwin J. Jones, Director, Virginia Cooperative Extension, Virginia Tech, Blacksburg; Jewel E. Hairston, Administrator, 1890 Extension Program, Virginia State, Petersburg.
May 1, 2009 |
Submitted by brad on Fri, 2009-06-12 13:49.
Our world has not rid itself of atrocity and genocide. What can modern high-tech do to help? In Bosnia, we used bombs. In Rwanda, we did next to nothing. In Darfur, very little. Here’s a proposal that seems expensive at first, but is in fact vastly cheaper than the military solutions people have either tried or been afraid to try. It’s the sunlight principle.
First, we would mass-produce a special video recording “phone” using the standard parts and tools of the cell phone industry. It would be small, light, and rechargeable from a car lighter plug, or possibly more slowly through a small solar cell on the back. It would cost a few hundred dollars to make, so that relief forces could airdrop tens or even hundreds of thousands of them over an area where atrocity is taking place. (If they are $400/pop, even 100,000 of them is 40 million dollars, a drop in the bucket compared to the cost of military operations.) They could also be smuggled in by relief workers on a smaller scale, or launched over borders in a pinch. Enough of them so that there are so many that anybody performing an atrocity will have to worry that there is a good chance that somebody hiding in bushes or in a house is recording it, and recording their face. This fear alone would reduce what took place.
Once the devices had recorded a video, they would need to upload it. It seems likely that in these situations the domestic cell system would not be available, or would be shut down to stop video uploads. However, that might not be true, and a version that uses existing cell systems might make sense, and be cheaper because the hardware is off the shelf. It is more likely that some other independent system would be used, based on the same technology but with slightly different protocols.
The anti-atrocity team would send aircraft over the area. These might be manned aircraft (presuming air superiority) or they might be very light, autonomous UAVs of the sort that already are getting cheap in price. These UAVs can be small, and not that high-powered, because they don’t need to do that much transmitting — just a beacon and a few commands and ACKs. The cameras on the ground will do the transmitting. In fact, the UAVs could quite possibly be balloons, again within the budget of aid organizations, not just nations. read more »
Submitted by brad on Sat, 2009-04-18 19:37.
My prior post about USB charging hubs in hotel rooms brought up the issue of security, as was the case for my hope for a world with bluetooth keyboards scattered around.
Is it possible to design our computers to let them connect to untrusted devices? Clearly to a degree, in that an ethernet connection is generally always untrusted. But USB was designed to be fully trusted, and that limits it.
Perhaps in the future, an OS can be designed to understand the difference between trusted and untrusted devices connected (wired or wirelessly) to a computer or phone. This might involve a different physical interface, or using the same physical interface, but a secure protocol by which devices can be identified (and then recognized when plugged in again) and tagged once as trusted the first time they are plugged in.
For example, an unknown keyboard is a risky thing to plug in. It could watch you type and remember passwords, or it could simply send fake keys to your computer to get it to install trojan software completely taking it over. But we might allow an untrusted keyboard to type plain text into our word processors or E-mail applications. However, we would have to switch to the trusted keyboard (which might just be a touch-screen keyboard on a phone or tablet) for anything dangerous, including of course entry of passwords, URLs and commands that go beyond text entry. Would this be tolerable, constantly switching like this, or would we just get used to it? We would want to mount the inferior keyboard very close to our comfy but untrusted one.
A mouse has the same issues. We might allow an untrusted mouse to move the pointer within a text entry window and to go to a set of menus that can’t do anything harmful on the machine, but would it drive us crazy to have to move to a different pointer to move out of the application? Alas, an untrusted mouse can (particularly if it waits until you are not looking) run applications, even bring up the on-screen keyboard most OSs have for the disabled, and then do anything with your computer.
It’s easier to trust output devices, like a printer. In fact, the main danger with plugging in an unknown USB printer is that a really nasty one might pretend to be a keyboard or CD-Rom to infect you. A peripheral bus that allows a device to only be an output device would be safer. Of course an untrusted printer could still record what you print.
An untrusted screen is a challenge. While mostly safe, one can imagine attacks. An untrusted screen might somehow get you to go to a special web-site. There, it might display something else, perhaps logins for a bank or other site so that it might capture the keys. Attacks here are difficult but not impossible, if I can control what you see. It might be important to have the trusted screen nearby somehow helping you to be sure the untrusted screen is being good. This is a much more involved attack than the simple attacks one can do by pretending to be a keyboard.
An untrusted disk (including a USB thumb drive) is actually today’s biggest risk. People pass around thumb drives all the time, and they can pretend to be auto-run CD-roms. In addition, we often copy files from them, and double click on files on them, which is risky. The OS should never allow code to auto-run from an untrusted disk, and should warn if files are double-clicked from them. Of course, even then you are not safe from traps inside the files themselves, even if the disk is just being a disk. Many companies try to establish very tight firewalls but it’s all for naught if they allow people to plug external drives and thumbsticks into the computers. Certain types of files (such as photos) are going to be safer than others (like executables and word processor files with macros or scripts.) Digital cameras, which often look like drives, are a must, and can probably be trusted to hand over jpegs and other image and video files.
A network connection is one of the things you can safely plug in. After all, a network connection should always be viewed as hostile, even one behind a firewall.
There is a risk in declaring a device trusted, for example, such as your home keyboard. It might be compromised later, and there is not much you can do about that. A common trick today is to install a key-logger in somebody’s keyboard to snoop on them. This is done not just by police but by suspicious spouses and corporate spies. Short of tamper-proof hardware and encryption, this is a difficult problem. For now, that’s too much cost to add to consumer devices.
Still, it sure would be nice to be able to go to a hotel and use their keyboard, mouse and monitor. It might be worth putting up with having to constantly switch back to get full sized input devices on computers that are trying to get smaller and smaller. But it would also require rewriting of a lot of software, since no program could be allowed to take input from an untrusted device unless it has been modified to understand such a protocol. For example, your e-mail program would need to be modified to declare that a text input box allows untrusted input. This gets harder in web browsing — each web page would need to have to declare, in its input boxes, whether untrusted input was allowed.
As a starter, however, the computer could come with a simple “clipboard editor” which brings up a box in which one can type and edit with untrusted input devices. Then, one could copy the edited text to the OS clipboard and, using the trusted mouse or keyboard, paste it into any application of choice. You could always get back to the special editing windows using the untrusted keyboard and mouse, you would have to use the trusted ones to leave that window. Cumbersome, but not as cumbersome as typing a long e-mail on an iPhone screen.
Submitted by brad on Thu, 2009-03-05 00:35.
I’m looking at you Ubuntu.
For some time now, the standard form for distributing a free OS (ie. Linux, *BSD) has been as a CD-ROM or DVD ISO file. You burn it to a CD, and you can boot and install from that, and also use the disk as a live CD.
There are a variety of pages with instructions on how to convert such an ISO into a bootable flash drive, and scripts and programs for linux and even for windows — for those installing linux on a windows box.
And these are great and I used one to make a bootable Ubuntu stick on my last install. And wow! It’s such a much nicer, faster experience compared to using CD that it’s silly to use CD on any system that can boot from a USB drive, and that’s most modern systems. With a zero seek time, it is much nicer.
So I now advocate going the other way. Give me a flash image I can dd to my flash drive, and a tool to turn that into an ISO if I need an ISO.
This has a number of useful advantages:
- I always want to try the live CD before installing, to make sure the hardware works in the new release. In fact, I even do that before upgrading most of the time.
- Of course, you don’t have old obsolete CDs lying around.
- Jumping to 1 gigabyte allows putting more on the distribution, including some important things that are missing these days, such as drivers and mdadm (the RAID control program.)
- Because flash is a dynamic medium, the install can be set up so that the user can, after copying the base distro, add files to the flash drive, such as important drivers — whatever they choose. An automatic script could even examine a machine and pull down new stuff that’s needed.
- You get a much faster and easier to use “rescue stick.”
- It’s easier to carry around.
- No need for an “alternate install” and perhaps easier as well to have the upgrader use the USB stick as a cache of packages during upgrades.
- At this point these things are really cheap. People give them away. You could sell them. This technique would also work for general external USB drives, or even plain old internal hard drives temporarily connected to a new machine being built if boot from USB is not practical. Great and really fast for eSata.
- Using filesystems designed not to wear out flash, the live stick can have a writable partition for /tmp, installed packages and modifications (with some security risk if you run untrusted code.)
Submitted by brad on Sat, 2009-02-14 19:34.
Product recalls have been around for a while. You get a notice in the mail. You either go into a dealer at some point, any point, for service, or you swap the product via the mail. Nicer recalls mail you a new product first and then you send in the old one, or sign a form saying you destroyed it. All well and good. Some recalls are done as “hidden warranties.” They are never announced, but if you go into the dealer with a problem they just fix it for free, long after the regular warranty, or fix it while working on something else. These usually are for items that don’t involve safety or high liability.
Today I had my first run-in with a recall of a connected electronic product. I purchased an “EyeFi” card for my sweetie for valentines day. This is an SD memory card with an wifi transmitter in it. You take pictures, and it stores them until it encounters a wifi network it knows. It then uploads the photos to your computer or to photo sharing sites. All sounds very nice.
When she put in the card and tried to initialize it, up popped a screen. “This card has a defect. Please give us your address and we’ll mail you a new one, and you can mail back the old one, and we’ll give you a credit in our store for your trouble.” All fine, but the product refused to let her register and use the product. We can’t even use the product for a few days to try it out (knowing it may lose photos.) What if I wanted to try it out to see if I was going to return it to the store. No luck. I could return it to the store as-is, but that’s work and may just get another one on the recall list.
This shows us the new dimension of the electronic recall. The product was remotely disabled to avoid liability for the company. We had no option to say, “Let us use the card until the new one arrives, we agree that it might fail or lose pictures.” For people who already had the card, I don’t know if it shut them down (possibly leaving them with no card) or let them continue with it. You have to agree on the form that you will not use the card any more.
This can really put a damper on a gift, when it refuses to even let you do a test the day you get it.
With electronic recall, all instances of a product can be shut down. This is similar to problems that people have had with automatic “upgrades” that actually remove features (like adding more DRM) or which fix you jailbreaking your iPhone. You don’t own the product any more. Companies are very worried about liability. They will “do the safe thing” which is shut their product down rather than let you take a risk. With other recalls, things happened on your schedule. You were even able to just decide not to do the recall. The company showed it had tried its best to convince you to do it, and could feel satisfied for having tried.
This is one of the risks I list in my essays on robocars. If a software flaw is found in a robocar (or any other product with physical risk) there will be pressure to “recall” the software and shut down people’s cars. Perhaps in extreme cases while they are driving on the street! The liability of being able to shut down the cars and not doing so once you are aware of a risk could result in huge punitive damages under the current legal system. So you play it safe.
But if people find their car shutting down because of some very slight risk, they will start wondering if they even want a car that can do that. Or even a memory card. Only with public pressure will we get the right to say, “I will take my own responsibility. You’ve informed me, I will decide when to take the product offline to get it fixed.”
Submitted by brad on Mon, 2008-09-29 22:40.
Most of us have had to stand in a long will-call line to pick up tickets. We probably even paid a ticket “service fee” for the privilege. Some places are helping by having online printable tickets with a bar code. However, that requires that they have networked bar code readers at the gate which can detect things like duplicate bar codes, and people seem to rather have giant lines and many staff rather than get such machines.
Can we do it better?
Well, for starters, it would be nice if tickets could be sent not as a printable bar code, but as a message to my cell phone. Perhaps a text message with coded string, which I could then display to a camera which does OCR of it. Same as a bar code, but I can actually get it while I am on the road and don’t have a printer. And I’m less likely to forget it.
Or let’s go a bit further and have a downloadable ticket application on the phone. The ticket application would use bluetooth and a deliberately short range reader. I would go up to the reader, and push a button on the cell phone, and it would talk over bluetooth with the ticket scanner and authenticate the use of my ticket. The scanner would then show a symbol or colour and my phone would show that symbol/colour to confirm to the gate staff that it was my phone that synced. (Otherwise it might have been the guy in line behind me.) The scanner would be just an ordinary laptop with bluetooth. You might be able to get away with just one (saving the need for networking) because it would be very fast. People would just walk by holding up their phones, and the gatekeeper would look at the screen of the laptop (hidden) and the screen of the phone, and as long as they matched wave through the number of people it shows on the laptop screen.
Alternately you could put the bluetooth antenna in a little faraday box to be sure it doesn’t talk to any other phone but the one in the box. Put phone in box, light goes on, take phone out and proceed.
One reason many will-calls are slow is they ask you to show ID, often your photo-ID or the credit card used to purchase the item. But here’s an interesting idea. When I purchase the ticket online, let me offer an image file with a photo. It could be my photo, or it could be the photo of the person I am buying the tickets for. It could be 3 photos if any one of those 3 people can pick up the ticket. You do not need to provide your real name, just the photo. The will call system would then inkjet print the photos on the outside of the envelope containing your tickets.
You do need some form of name or code, so the agent can find the envelope, or type the name in the computer to see the records. When the agent gets the envelope, identification will be easy. Look at the photo on the envelope, and see if it’s the person at the ticket window. If so, hand it over, and you’re done! No need to get out cards or hand them back and forth.
A great company to implement this would be paypal. I could pay with paypal, not revealing my name (just an E-mail address) and paypal could have a photo stored, and forward it on to the ticket seller if I check the box to do this. The ticket seller never knows my name, just my picture. You may think it’s scary for people to get your picture, but in fact it’s scarier to give them your name. They can collect and share data with you under your name. Your picture is not very useful for this, at least not yet, and if you like you can use one of many different pictures each time — you can’t keep using different names if you need to show ID.
This could still be done with credit cards. Many credit cards offer a “virtual credit card number” system which will generate one-time card numbers for online transactions. They could set these up so you don’t have to offer a real name or address, just the photo. When picking up the item, all you need is your face.
This doesn’t work if it’s an over-21 venue, alas. They still want photo ID, but they only need to look at it, they don’t have to record the name.
It would be more interesting if one could design a system so that people can find their own ticket envelopes. The guard would let you into the room with the ticket envelopes, and let you find yours, and then you can leave by showing your face is on the envelope. The problem is, what if you also palmed somebody else’s envelope and then claimed yours, or said you couldn’t find yours? That needs a pretty watchful guard which doesn’t really save on staff as we’re hoping. It might be possible to have the tickets in a series of closed boxes. You know your box number (it was given to you, or you selected it in advance) so you get your box and bring it to the gate person, who opens it and pulls out your ticket for you, confirming your face. Then the box is closed and returned. Make opening the boxes very noisy.
I also thought that for Burning Man, which apparently had a will-call problem this year, you could just require all people fetching their ticket be naked. For those not willing, they could do regular will-call where the ticket agent finds the envelope. :-)
I’ve noted before that, absent the need of the TSA to know all our names, this is how boarding passes should work. You buy a ticket, provide a photo of the person who is to fly, and the gate agent just looks to see if the face on the screen is the person flying, no need to get out ID, or tell the airline your name.
Submitted by brad on Tue, 2008-05-27 20:49.
Hard disks fail. If you prepared properly, you have a backup, or you swap out disks when they first start reporting problems. If you prepare really well you have offsite backup (which is getting easier and easier to do over the internet.)
One way to protect yourself from disk failures is RAID, especially RAID-5. With RAID, several disks act together as one. The simplest protecting RAID, RAID-1, just has 2 disks which work in parallel, known as mirroring. Everything you write is copied to both. If one fails, you still have the other, with all your data. It’s good, but twice as expensive.
RAID-5 is cleverer. It uses 3 or more disks, and uses error correction techniques so that you can store, for example, 2 disks worth of data on 3 disks. So it’s only 50% more expensive. RAID-5 can be done with many more disks — for example with 5 disks you get 4 disks worth of data, and it’s only 25% more expensive. However, having 5 disks is beyond most systems and has its own secret risk — if 2 of the 5 disks fail at once — and this does happen — you lose all 4 disks worth of data, not just 2 disks worth. (RAID-6 for really large arrays of disks, survives 2 failures but not 3.)
Now most people who put in RAID do it for more than data protection. After all, good sysadmins are doing regular backups. They do it because with RAID, the computer doesn’t even stop when a disk fails. You connect up a new disk live to the computer (which you can do with some systems) and it is recreated from the working disks, and you never miss a beat. This is pretty important with a major server.
But RAID has value to those who are not in the 99.99% uptime community. Those who are not good at doing manual backups, but who want to be protected from the inevitable disk failures. Today it is hard to set up, or expensive, or both. There are some external boxes like the “readynas” that make it reasonably easy for external disks, but they don’t have the bandwidth to be your full time disks.
RAID-5 on old IDE systems was hard, they usually could truly talk to only 2 disks at a time. The new SATA bus is much better, as many motherboards have 4 connectors, though soon one will be required by blu-ray drives. read more »
Submitted by brad on Thu, 2008-05-15 13:56.
Recently we at the EFF have been trying to fight new rulings about the power of U.S. customs. Right now, it’s been ruled they can search your laptop, taking a complete copy of your drive, even if they don’t have the normally required reasons to suspect you of a crime. The simple fact that you’re crossing the border gives them extraordinary power.
We would like to see that changed, but until then what can be done? You can use various software to encrypt your hard drive — there are free packages like truecrypt, and many laptops come with this as an option — but most people find having to enter a password every time you boot to be a pain. And customs can threaten to detain you until you give them the password.
There are some tricks you can pull, like having a special inner-drive with a second password that they don’t even know to ask about. You can put your most private data there. But again, people don’t use systems with complex UIs unless they feel really motivated.
What we need is a system that is effectively transparent most of the time. However, you could take special actions when going through customs or otherwise having your laptop be out of your control. read more »
Submitted by brad on Sat, 2008-05-10 18:46.
It seems that half the programs I try and install under Windows want to have a “daemon” process with them, which is to say a portion of the program that is always running and which gets a little task-tray icon from which it can be controlled. Usually they want to also be run at boot time. In Windows parlance this is called a service.
There are too many of them, and they don’t all need to be there. Microsoft noticed this, and started having Windows detect if task tray icons were too static. If they are it hides them. This doesn’t work very well — they even hide their own icon for removing hardware, which of course is going to be static most of the time. And of course some programs now play games to make their icons appear non-static so they will stay visible. A pointless arms race.
All these daemons eat up memory, and some of them eat up CPU. They tend to slow the boot of the machine too. And usually not to do very much — mostly to wait for some event, like being clicked, or hardware being plugged in, or an OS/internet event. And the worst of them on their menu don’t even have a way to shut them down.
I would like to see the creation of a master deaemon/service program. This program would be running all the time, and it would provide a basic scripting language to perform daemon functions. Programs that just need a simple daemon, with a menu or waiting for events, would be strongly encouraged to prepare it in this scripting language, and install it through the master daemon. That way they take up a few kilobytes, not megabytes, and don’t take long to load. The scripting language should be able to react at least in a basic way to all the OS hooks, events and callbacks. It need not do much with them — mainly it would run a real module of the program that would have had a daemon. If the events are fast and furious and don’t pause, this program could stay resident and become a real daemon.
But having a stand alone program would be discouraged, certainly for boring purposes like checking for updates, overseeing other programs and waiting for events. The master program itself could get regular updates, as features are added to it as needed by would-be daemons.
Unix started with this philosophy. Most internet servers are started up by inetd, which listens on all the server ports you tell it, and fires up a server if somebody tries to connect. Only programs with very frequent requests, like E-mail and web serving, are supposed to keep something constantly running.
The problem is, every software package is convinced it’s the most important program on the system, and that the user mostly runs nothing but that program. So they act like they own the place. We need a way to only let them do that if they truly need it.
Submitted by brad on Fri, 2008-05-09 00:14.
I’m scanning my documents on an ADF document scanner now, and it’s largely pretty impressive, but I’m surprised at some things the system won’t do.
Double page feeding is the bane of document scanning. To prevent it, many scanners offer methods of double feed detection, including ultrasonic detection of double thickness and detection when one page is suddenly longer than all the others (because it’s really two.)
There are a number of other tricks they could do, I think. I think a paper feeder that used air suction or gecko-foot van-der-waals force pluckers on both sides of a page to try to pull the sides in two different directions could help not just detect, but eliminate such feeds.
However, the most the double feed detectors do is signal an exception to stop the scan. Which means work re-feeding and a need to stand by.
However, many documents have page numbers. And we’re going to OCR them and the OCR engine is pretty good at detecting page numbers (mostly out of desire to remove them.) However, it seems to me a good approach would be to look for gaps in the page numbers, especially combined with the other results of a double feed. Then don’t stop the scan, just keep going, and report to the operator which pages need to be scanned again. Those would be scanned, their number extracted, and they would be inserted in the right place in the final document.
Of course, it’s not perfect. Sometimes page numbers are not put on blank pages, and some documents number only within chapters. So you might not catch everything, but you could catch a lot of stuff. Operators could quickly discern the page numbering scheme (though I think the OCR could do this too) to guide the effort.
I’m seeking a maximum convenience workflow. I think to do that the best plan is to have several scanners going, and the OCR after the fact in the background. That way there’s always something for the operator to do — fixing bad feeds, loading new documents, naming them — for maximum throughput. Though I also would hope the OCR software could do better at naming the documents for you, or at least suggesting names. Perhaps it can, the manual for Omnipage is pretty sparse.
While some higher end scanners do have the scanner figure out the size of the page (at least the length) I am not sure why it isn’t a trivial feature for all ADF scanners to do this. My $100 Strobe sheetfed scanner does it. That my $6,000 (retail) FI-5650 needs extra software seems odd to me.
Submitted by brad on Tue, 2008-05-06 16:25.
PCs can go into standby mode (just enough power to preserve the RAM and do wake-on-lan) and into hibernate mode (where they write out the RAM to disk, shut down entirely and restore from disk later) as well as fully shut down.
Standby mode comes back up very fast, and should be routinely used on desktops. In fact, non-server PCs should consider doing it as a sort of screen saver since the restart can be so quick. It’s also popular on laptops but does drain the battery in a few days keeping the RAM alive. Many laptops will wake up briefly to hibernate if left in standby so long that the battery gets low, which is good.
How about this option: Write the ram contents out to disk, but also keep the ram alive. When the user wants to restart, they can restart instantly, unless something happened to the ram. If there was a power flicker or other trouble, notice the ram is bad and restart from disk. Usually you don’t care too much about the extra time needed to write out to disk when suspending, other than for psychological reasons where you want to be really sure the computer is off before leaving it. It’s when you come back to the computer that you want instant-on.
In fact, since RAM doesn’t actually fail all that quickly, you might even find you can restore from RAM after a brief power flicker. In that case, you would want to store a checksum for all blocks of RAM, and restore any from disk that don’t match the checksum.
To go further, one could also hibernate to newer generations of fast flash memory. Flash memory is getting quite cheap, and while older generations aren’t that quick, they seek instantaneously. This allows you to reboot a machine with its memory “paged out” to flash, and swap in pages at random as they are needed. This would allow a special sort of hybrid restore:
- Predict in advance which pages are highly used, and which are enough to get the most basic functions of the OS up. Write them out to a special contiguous block of hibernation disk. Then write out the rest, to disk and flash.
- When turning on again, read this block of contiguous disk and go “live.” Any pages needed can then be paged in from the flash memory as needed, or if the flash wasn’t big enough, unlikely pages can come from disk.
- In the background, restore the rest of the pages from the faster disk. Eventually you are fully back to ram.
This would allow users to get a fairly fast restore, even from full-off hibernation. If they click on a rarely used program that was in ram, it might be slow as stuff pages in, but still not as bad as waiting for the whole restore.
Submitted by brad on Thu, 2008-02-21 12:44.
A big trend in systems operation these days is the use of virtual machines — software systems which emulate a standalone machine so you can run a guest operating system as a program on top of another (host) OS. This has become particularly popular for companies selling web hosting. They take one fast machine and run many VMs on it, so that each customer has the illusion of a standalone machine, on which they can do anything. It’s also used for security testing and honeypots.
The virtual hosting is great. Typical web activity is “bursty.” You would like to run at a low level most of the time, but occasionally burst to higher capacity. A good VM environment will do that well. A dedicated machine has you pay for full capacity all the time when you only need it rarely. Cloud computing goes beyond this.
However, the main limit to a virtual machine’s capacity is memory. Virtual host vendors price their machines mostly on how much RAM they get. And a virtual host with twice the RAM often costs twice as much. This is all based on the machine’s physical ram. A typical vendor might take a machine with 4gb, keep 256mb for the host and then sell 15 virtual machines with 256mb of ram. They will also let you “burst” your ram, either into spare capacity or into what the other customers are not using at the time, but if you do this for too long they will just randomly kill processes on your machine, so you don’t want to depend on this.
The problem is when they give you 256MB of ram, that’s what you get. A dedicated linux server with 256mb of ram will actually run fairly well, because it uses paging to disk. The server loads many programs, but a lot of the memory used for these programs (particularly the code) is used rarely, if ever, and swaps out to disk. So your 256mb holds the most important pages of ram. If you have more than 256mb of important, regularly used ram, you’ll thrash (but not die) and know you need to buy more.
The virtual machines, however, don’t give you swap space. Everything stays in ram. And the host doesn’t swap it either, because that would not be fair. If one VM were regularly swapping to disk, this would slow the whole system down for everybody. One could build a fair allocation for that but I have not heard of it.
In addition, another big memory saving is lost — shared memory. In a typical system, when two processes use the same shared library or same program, this is loaded into memory only once. It’s read-only so you don’t need to have two copies. But on a big virtual machine, we have 15 copies of all the standard stuff — 15 kernels, 15 MYSQL servers, 15 web servers, 15 of just about everything. It’s very wasteful.
So I wonder if it might be possible to do one of the following:
- Design the VM so that all binaries and shared libraries can be mounted from a special read-only filesystem which is actually on the host. This would be an overlay filesystem so that individual virtual machines could change it if need be. The guest kernel, however, would be able to load pages from these files, and they would be shared with any other virtual machine loading the same file.
- Write a daemon that regularly uses spare CPU to scan the pages of each virtual machine, hashing them. When two pages turn out to be identical, release one and have both VMs use the common copy. Mark it so that if one writes to it, a duplicate is created again. When new programs start it would take extra RAM, but within a few minutes the memory would be shared.
These techniques require either a very clever virtualizer or modified guests, but their savings are so worthwhile that everybody would want to do it this way on any highly loaded virtual machine. Of course, that goes against the concept of “run anything you like” and makes it “run what you like, but certain standard systems are much cheaper.”
This, and allowing some form of fair swapping, could cause a serious increase in the performance and cost of VMs.
Submitted by brad on Tue, 2008-02-19 21:11.
If you have read my articles on power you know I yearn for the days when we get smart power so we have have universal supplies that power everything. This hit home when we got a new Thinkpad Z61 model, which uses a new power adapter which provides 20 volts at 4.5 amps and uses a new, quite rare power tip which is 8mm in diameter. For almost a decade, thinkpads used 16.5 volts and used a fairly standard 5.5mm plug. It go so that some companies standardized on Thinkpads and put cheap 16 volt TP power supplies in all the conference rooms, allowing employees to just bring their laptops in with no hassle.
Lenovo pissed off their customers with this move. I have perhaps 5 older power supplies, including one each at two desks, one that stays in the laptop bag for travel, one downstairs and one running an older ThinkPad. They are no good to me on the new computer.
Lenovo says they knew this would annoy people, and did it because they needed more power in their laptops, but could not increase the current in the older plug. I’m not quite sure why they need more power — the newer processors are actually lower wattage — but they did.
Here’s something they could have done to make it better. read more »
Submitted by brad on Sat, 2008-01-12 16:33.
I’ve written before about both the desire for universal dc power and more simply universal laptop power at meeting room desks.
Today I want to report we’re getting a lot closer. A new generation of cheap “buck and boost” ICs which can handle more serious wattages with good efficiency has come to the market. This means cheap DC to DC conversion, both increasing and decreasing voltages. More and more equipment is now able to take a serious range of input voltages, and also to generate them. Being able to use any voltage is important for battery powered devices, since batteries start out with a high voltage (higher than the one they are rated for) and drop over their time to around 2/3s of that before they are viewed as depleted. (With some batteries, heavy depletion can really hurt their life. Some are more able to handle it.)
With a simple buck converter chip, at a cost of about 10-15% of the energy, you get a constant voltage out to matter what the battery is putting out. This means more reliable power and also the ability to use the full capacity of the battery, if you need it and it won’t cause too much damage. These same chips are in universal laptop supplies. Most of these supplies use special magic tips which fit the device they are powering and also tell the supply what voltage and current it needs. read more »
Submitted by brad on Tue, 2007-11-13 13:20.
Ok, I haven't had a new laptop in a while so perhaps this already happens, but I'm now carrying more devices that can charge off the USB power, including my cell phone. It's only 2.5 watts, but it's good enough for many purposes.
However, my laptops, and desktops, do not provide USB power when in standby or off. So how about a physical or soft switch to enable that? Or even a smart mode in the US that lets you list what devices you want to keep powered and which ones you don't? (This would probably keep all devices powered if any one such device is connected, unless you had individual power control for each plug.)
This would only be when on AC power of course, not on battery unless explicitly asked for as an emergency need.
To get really smart a protocol could be developed where the computer can ask the USB device if it needs power. A fully charged device that plans to sleep would say no. A device needing charge could say yes.
Of course, you only want to do this if the power supply can efficiently generate 5 volts. Some PC power supplies are not efficient at low loads and so may not be a good choice for this, and smaller power supplies should be used.
Submitted by brad on Tue, 2007-07-10 00:42.
For much of history, we’ve used removable media for backup. We’ve used tapes of various types, floppy disks, disk cartridges, and burnable optical disks. We take the removable media and keep a copy offsite if we’re good, but otherwise they sit for a few decades until they can’t be read, either because they degraded or we can’t find a reader for the medium any more.
But I now declare this era over. Disk drives are so cheap — 25 cents/gb and falling, that it no longer makes sense to do backups to anything but hard disks. We may use external USB drives that are removable, but at this point our backups are not offline, they are online. Thanks to the internet, I even do offsite backup to live storage. I sync up over the internet at night, and if I get too many changes (like after an OS install, or a new crop of photos) I write the changes to a removable hard disk and carry it over to the offsite hard disk.
Of course, these hard drives will fail, perhaps even faster than CD-roms or floppies. But the key factor is that the storage is online rather than offline, and each new disk is 2 to 3 times larger than the one it replaced. What this means is that as we change out our disks, we just copy our old online archives to our new online disk. By constantly moving the data to newer and newer media — and storing it redundantly with online, offsite backup, the data are protected from the death that removable media eventually suffer. So long as disks keep getting bigger and cheaper, we won’t lose anything, except by beng lazy. And soon, our systems will get more automated at this, so it’s hard to set up a computer that isn’t backed up online and remotely. We may still lose things because we lose encryption keys, but it won’t be for media.
Thus, oddly, the period of the latter part of the 20th century will be a sort of “dark ages” to future data archaeologists. Those disks will be lost. The media may be around, but you will have to do a lot of work to recover them — manual work. However, data from the early 21st onward will be there unless it was actively deleted or encrypted.
Of course this has good and bad consequences. Good for historians. Perhaps not so good for privacy.
Submitted by brad on Tue, 2007-07-03 15:15.
Hotels are now commonly sporting flat widescreen TVs, usually LCD HDTVs at the 720p resolution, which is 1280 x 720 or similar. Some of these TVs have VGA ports or HDMI (DVI) ports, or they have HDTV analog component video (which is found on some laptops but not too many.) While 720p resolution is not as good as the screens on many laptops, it makes a world of difference on a PDA. As our phone/PDA devices become more like the iPhone, it would be very interesting to see hotels guarantee that their room offers the combination of:
- A bluetooth keyboard (with USB and mini-USB as a backup)
- A similar optical mouse
- A means to get video into the HDTV
- Of course, wireless internet
- Our dreamed of universal DC power jack (or possibly inductive charging.)
Tiny devices like the iPhone won’t sport VGA or even component video out 7 pin connectors, though they might do HDMI. It’s also not out of the question to go a step further and do a remote screen protocol like VNC over the wireless ethernet or bluetooth.
This would engender a world where you carry a tiny device like the iPhone, which is all touchscreen for when you are using it in the mobile environment. However, when you sit down in your hotel room (or a few other places) you could use it like a full computer with a full screen and keyboard. (There are also quite compact real-key bluetooth keyboards and mice which travelers could also bring. Indeed, since the iPhone depends on a multitouch interface, an ordinary mouse might not be enough for it, but you could always use its screen for such pointing, effectively using the device as the touchpad.)
Such stations need not simply be in hotels. Smaller displays (which are now quite cheap) could also be present at workstations on conference tables or meeting rooms, or even for rent in public. Of course rental PCs in public are very common at internet cafes and airport kiosks, but using our own device is more tuned to our needs and more secure (though using a rented keyboard presents security risks.)
One could even imagine stations like these randomly scattered around cities behind walls. Many retailers today are putting HDTV flat panels in their windows instead of signs, and this will become a more popular trend. Imagine being able to borrow (for free or for a rental fee) such screens for a short time to do a serious round of web surfing on your portable device with high resolution, and local wifi bandwidth. Such a screen could not provide you with a keyboard or mouse easily, but the surfing experience would be much better than the typical mobile device surfing experience, even the iPhone model of seeing a blurry, full-size web page and using multitouch to zoom in on the relevant parts. Using a protocol like vnc could provide a good surfing experience for pedestrians.
Cars are also more commonly becoming equipped with screens, and they are another place we like to do mobile surfing. While the car’s computer should let you surf directly, there is merit in being able to use that screen as a temporary large screen for one’s mobile device.
Until we either get really good VR glasses or bright tiny projectors, screen size is going to be an issue in mobile devices. A world full of larger screens that can be grabbed for a few minutes use may be a good answer.
Submitted by brad on Fri, 2007-06-08 14:43.
For many of us, E-mail has become our most fundamental tool. It is not just the way we communicate with friends and colleagues, it is the way that a large chunk of the tasks on our “to do” lists and calendars arrive. Of course, many E-mail programs like Outlook come integrated with a calendar program and a to-do list, but the integration is marginal at best. (Integration with the contact manager/address book is usually the top priority.)
If you’re like me you have a nasty habit. You leave messages in your inbox that you need to deal with if you can’t resolve them with a quick reply when you read them. And then those messages often drift down in the box, off the first screen. As a result, they are dealt with much later or not at all. With luck the person mails you again to remind you of the pending task.
There are many time management systems and philosophies out there, of course. A common theme is to manage your to-do list and calendar well, and to understand what you will do and not do, and when you will do it if not right away.
I think it’s time to integrate our time management concepts with our E-mail. To realize that a large number of emails or threads are also a task, and should be bound together with the time manager’s concept of a task.
For example, one way to “file” an E-mail would be to the calendar or a day oriented to-do list. You might take an E-mail and say, “I need 20 minutes to do this by Friday” or “I’ll do this after my meeting with the boss tomorrow.” The task would be tied to the E-mail. Most often, the tasks would not be tied to a specific time the way calendar entries are, but would just be given a rough block of time within a rough window of hours or days.
It would be useful to add these “when to do it” attributes to E-mails, because now delegating a task to somebody else can be as simple as forwarding the E-mail-message-as-task to them.
In fact, because, as I have noted, I like calendars with free-form input (ie. saying “Lunch with Peter 1pm tomorrow” and having the calender understand exactly what to do with it) it makes sense to consider the E-mail window as a primary means of input to the calendar. For example, one might add calendar entries by emailing them to a special address that is processed by the calendar. (That’s a useful idea for any calendar, even one not tied at all to the E-mail program.)
One should also be able to assign tasks to places (a concept from the “Getting Things Done” book I have had recommended to me.) In this case, items that will be done when one is shopping, or going out to a specific meeting, could be synced or sent appropriately to one’s mobile device, but all with the E-mail metaphor.
Because there are different philosophies of time management, all with their fans, one monolithic e-mail/time/calendar/todo program may not be the perfect answer. A plug-in architecture that lets time managers integrate nicely with E-mail could be a better way to do it.
Some of these concepts apply to the shared calendar concepts I wrote about last month.
Submitted by brad on Mon, 2007-06-04 11:01.
Here’s a new approach to linux adoption. Create a linux distro which converts a Windows machine to linux, marketed as a way to solve many of your virus/malware/phishing woes.
Yes, for a long time linux distros have installed themselves on top of a windows machine dual-boot. And there are distros that can run in a VM on windows, or look windows like, but here’s a set of steps to go much further, thanks to how cheap disk space is today. read more »
- Yes, the distro keeps the Windows install around dual boot, but it also builds a virtual machine so it can be run under linux. Of course hardware drivers differ when running under a VM, so this is non-trivial, and Windows XP and later will claim they are stolen if they wake up in different hardware. You may have to call Microsoft, which they may eventually try to stop.
- Look through the Windows copy and see what apps are installed. For apps that migrate well to linux, either because they have equivalents or run at silver or gold level under Wine, move them into linux. Extract their settings and files and move those into the linux environment. Of course this is easiest to do when you have something like Firefox as the browser, but IE settings and bookmarks can also be imported.
- Examine the windows registry for other OS settings, desktop behaviours etc. Import them into a windows-like linux desktop. Ideally when it boots up, the user will see it looking and feeling a lot like their windows environment.
- Using remote window protocols, it’s possible to run windows programs in a virtual machine with their window on the X desktop. Try this for some apps, though understand some things like inter-program communication may not do as well.
- Next, offer programs directly in the virtual machine as another desktop. Put the windows programs on the windows-like “start” menu, but have them fire up the program in the virtual machine, or possibly even fire up the VM as needed. Again, memory is getting very cheap.
- Strongly encourage the Windows VM be operated in a checkpointing manner, where it is regularly reverted to a base state, if this is possible.
- The linux box, sitting outside the windows VM, can examine its TCP traffic to check for possible infections or strange traffic to unusual sites. A database like the siteadvisor one can help spot these unusual things, and encourage restoring the windows box back to a safe checkpoint.
Submitted by brad on Sun, 2007-04-15 16:45.
The use of virtual machines is getting very popular in the web hosting world. Particularly exciting to many people is Amazon.com’s EC2 — which means Elastic Compute Cloud. It’s a large pool of virtual machines that you can rent by the hour. I know people planning on basing whole companies on this system, because they can build an application that scales up by adding more virtual machines on demand. It’s decently priced and a lot cheaper than building it yourself in most cases.
In many ways, something like EC2 would be great for all those web sites which deal with the “slashdot” effect. I hope to see web hosters, servers and web applications just naturally allow scaling through the addition of extra machines. This typically means either some round-robin-DNS, or a master server that does redirects to a pool of servers, or a master cache that processes the data from a pool of servers, or a few other methods. Dealing with persistent state that can’t be kept in cookies requires a shared database among all the servers, which may make the database the limiting factor. Rumours suggest Amazon will release an SQL interface to their internal storage system which presumably is highly scalable, solving that problem.
As noted, this would be great for small to medium web sites. They can mostly run on a single server, but if they ever see a giant burst of traffic, for example by being linked to from a highly popular site, they can in minutes bring up extra servers to share the load. I’ve suggested this approach for the Battlestar Galactica Wiki I’ve been using — normally their load is modest, but while the show is on, each week, predictably, they get such a huge load of traffic when the show actually airs that they have to lock the wiki down. They have tried to solve this the old fashioned way — buying bigger servers — but that’s a waste when they really just need one day a week, 22 weeks a year, of high capacity.
However, I digress. What I really want to talk about is using such systems to get access to all sorts of platforms. As I’ve noted before, linux is a huge mishmash of platforms. There are many revisions of Ubuntu, Fedora, SuSE, Debian, Gentoo and many others out there. Not just the current release, but all the past releases, in both stable, testing and unstable branches. On top of that there are many versions of the BSD variants. read more »
Submitted by brad on Sun, 2007-03-04 19:50.
Most of us, when we travel, put appointments we will have while on the road into our calendars. And we usually enter them in local time. ie. if I have a 1pm appointment in New York, I set it for 1pm not 10am in my Pacific home time zone. While some calendar programs let you specify the time zone for an event, most people don't, and many people also don't change the time zone when they cross a border, at least not right away. (I presume that some cell phone PDAs pick up the new time from the cell network and import it into the PDA, if the network provides that.) Many PDAs don't really even let you set the time zone, just the time.
Here's an idea that's simple for the user. Most people put their flights into their calendars. In fact, most of the airline web sites now let you download your flight details right into your calendar. Those flight details include flight times and the airport codes.
So the calendar software should notice the flight, look up the destination airport code, and trigger a time zone change during the flight. This would also let the flight duration look correct in the calendar view window, though it would mean some "days" would be longer than others, and hours would repeat or be missing in the display.
You could also manually enter magic entries like "TZ to PST" or similar which the calendar could understand as a command to change the zone at that time.
Of course, I could go on many long rants about the things lacking from current calendar software, and perhaps at some point I will, but this one struck me as interesting because, in the downloaded case, the UI for the user is close to invisible, and I always like that.
It becomes important when we start importing our "presence" from our calendar, or get alerts from our devices about events, we don't want these things to trigger in the wrong time zone. |
AN EDUCATOR'S GUIDE:
MEETING THE NEEDS OF THE EHLERS-DANLOS CHILD
A PARENT'S GUIDE:
HELPING YOUR CHILD SUCCEED AT SCHOOL
The material in this guide has been prepared to help educators better understand and provide for the needs of the Ehlers-Danlos student. Each person with Ehlers-Danlos Syndrome is affected differently; therefore, the needs will be different for every student. The information has been compiled from a variety of resources including parents,administrators, teachers, guidance counselors and special education teachers. My sincere thanks to those who helped in this effort by sharing and providing personal and professional information.
|Peggy Rocha Snuggs |
|Ehlers-Danlos National Foundation |
|November, 2003 |
Before leaving the Board at the end of 2004, Peggy Rocha Snuggs was Director at Large for Children's Projects & Education, Board of Directors, Ehlers-Danlos National Foundation. She is an educator in the public school system in Tampa, Florida, and the mother of an EDS child.
Many teachers and schools have contacted the Ehlers-Danlos National Foundation with questions about EDS, ranging from basic concerns for the safety of the student, to questions about accommodations to meet the needs of the EDS child. This booklet was prepared to answer some of the questions asked most commonly by school personnel. Additional information on medical aspects of this syndrome is available from the Ehlers-Danlos National Foundation, or at www.ednf.org
Introduction to EDS
What is Ehlers-Danlos Syndrome (EDS)
The Ehlers-Danlos Syndrome (EDS) is a heterogeneous group of heritable connective tissue disorders affecting approximately 1 in 5,000 to 10,000 men and women of all ethnic backgrounds. EDS is named for two physicians (Ehlers and Danlos) who described the forms of the condition in the early 1900s. At least six forms of Ehlers-Danlos Syndrome have been described, which are not graduations in severity, but represent distinct disorders which "run-true" in a family. EDS is characterized by hyperextenisve skin, hypermobile joints, easy bruisability of the skin, and a bleeding diathesis. There are six major types of EDS, which are classified according to their different manifestation of signs and symptoms. Individuals with EDS have a defect in their connective tissue, the tissue which provides support to many parts of the body such as the skin, muscles and ligaments. The fragile skin and unstable joints found in EDS are the result of faulty collagen, the protein which acts as a "glue" in the body, adding strength and elasticity to connective tissue.
What medical problems are associated with EDS?
Due to the different types of EDS, and the varying degree of severity between individuals with EDS, it is difficult to generalize. However, most often, the manifestations are skin and joint related and may include:
Soft velvet-like skin; variable hyperextensibility (stretchiness); fragile skin that tears or bruises easily (bruising may be severe); severe scarring; slow and poor wound healing; and fleshy lesions (molluscoid pseudotumors) associated with scars over pressure areas.
Joint hypermobility; loose/unstable joints which are prone to frequent dislocations and/or subluxions; joint pain; hyperextensive joints (they move beyond the joint's normal range); early onset of osteoarthritis.
Chronic, early onset, debilitating musculoskeletal pain; arterial/intestinal/uterine fragility or rupture (usually associated with Vascular Type of EDS); scoliosis; poor muscle tone; mitral valve prolapse; gum disease and vision problems.
Medical Emergencies at School
What medical emergencies may arise?
It is wise to make all school staff, including lunchroom supervisors, nurses, and office, aware that a medical emergency may occur and to outline a basic plan of action. It is important to have a prepared, parent approved plan ahead on file. Current and updated parent contact information should be available at all times.
These often occur due to the loose ligaments and commonly affect the knees, shoulders and hips, but may affect the fingers and wrists or other joints as well. Activity may cause the joint to lock or overextend.
A physician may prescribe bracing to stabilize joints, or surgical repair of joints. Ice packs should be kept on hand. Occupational therapists and or physical therapists can help students learn how to strengthen their muscles and teach them how to properly use and preserve their joints.
Cuts/lacerations may occur and may vary from minor to major wounds. Gaping wounds should be handled with care. Scarring is not uncommon. Proper repair of these wounds by a physician may be necessary. Excessive bleeding may occur with any cut. Severe bruising is also common.
Excessive sun exposure should be avoided and sunscreen is recommended for outdoor activities.
Vascular Type of EDS:
The EDNF has prepared a CD-ROM and 20 page handbook to educate emergency room physicians and other health care professionals about the vascular type of EDS. This CD-Rom contains suggestions on life-saving surgical and postoperative techniques. In a trauma situation, time is of the essence. It would be beneficial for a school to retain a copy to accompany the child to the emergency room if it should become necessary. The CD-ROM is available from the EDNF shopping cart. To acquire a copy, please go to Shop at EDNF.
Some children with EDS may complain of severe gastro pain and frequent stomach aches. A family plan, as well as a school plan for dealing with chronic stomach pain should be addressed with the parent and the teachers.
STEPS FOR ACADEMIC SUCCESS
Step 1: Providing a Safe Classroom/Physical Environment
Although each person is affected differently by Ehlers-Danlos syndrome, the need for a safe environment at school is important to avoid injury.
When possible, carpeted areas are preferable to slippery floors, and avoiding the use of stairs in favor of elevators. An evaluation of appropriate seating by a school based occupational therapist should be considered.
Step 2: Providing Appropriate Physical Education - Sports, and Elective Classes
It is important, for emotional well-being that students with EDS try to lead as normal a life as possible. This includes playing with other children and participating as much as possible in physical activity. Students with EDS will have different needs or restrictions concerning physical education and sports. Where there is only mild hypermobility and the skin is not fragile, physical education and sports may be allowed with an adapted program. When guiding students to his/her choices of electives and activities, it is helpful to consider whether the activity could cause stress to joints or possible injury. Common sense should prevail when designing sport programs. Contact sports or activities that require vigorous exertion, heavy lifting, blows to the head or chest or excessive strain on joints and ligaments may not be the best choice. Sports such as swimming & golf may be a more logical alternative if done at an easy pace with opportunity for resting.
Certain musical instruments may stress some joints, where others may be more acceptable. Each activity should be accessed based on individual student abilities and restrictions.
Most students with EDS do not want to appear different. A student may experience all or some if the following which would affect their success in physical activities: impaired mobility, weak hand control, and/or poor coordination. Being forced to participate in physical education classes may cause the EDS child to become disheartened and embarrassed. Asking the child to sit on the sidelines and watch may cause the child to feel isolated and different. Ideally, alternative classes should be available if physical education classes are not recommended. Some schools allow health/related classes to be substituted for physical education credits, and some waiver the credit altogether in favor of a more suitable elective.
Step 3: Cognitive Profiles and Providing Testing Accommodations
There is no evidence to indicate that EDS in itself, causes learning difficulties. Premature birth is a complication associated with EDS and these children may experience the delays often associated with it. There are some speech and hearing and visual problems associated with EDS that medical specialist in the field can diagnose.
Formal IQ and performance testing can provide useful information. However, careful interpretation of formal evaluations should be made to assure that any occurrence of verbal-performance discrepancy, with performance scores being lower than verbal scores, is not due to motor problems instead of learning disability.
Teachers should be aware of potential difficulties due to hand-wrist hypermobility which are associated with reduced scores on performance tests. Students may need extra time or may require alternate forms of assess-ing performance that requires less motor input. Wrist or finger splints, large or padded pens or pencils or pencil grips may be helpful. Some students may need additional testing accommodations such as oral evaluations.
Step 4: Consideration for Chronic Absences
Due to frequent absences, the EDS child may need additional time, additional tutorials, and other child-specific accommodations. In some cases, at-home instructions may need to be provided for long term recoveries.
Allowing special pupil assignment to a school closer to parent's home or work is a viable consideration.
Some schools post homework and class work assignments on a website for students who are absent. If the school has such a site, parents should be made aware of it.
Some schools offer incentive grade points or exam exemptions based on attendance. The EDS child should not be penalized for EDS related absences.
Step 5: Providing Teacher Support in the Classroom
The inability to participate in some peer activities, the need for special accommodations, and the sense of "feeling different" may lead to frustration and isolation. Helping the EDS child feel accepted is something a teacher can help with. Addressing these issues by frequent discussions with the child about his or her feelings, identifying and encouraging development of other talents, and inclusion in activities that are not restricted are positive ways to help. Many EDS children feel comfortable talking about EDS, and may be willing to discuss or report on EDS in class. Some however are embarrassed and choose not to discuss their condition.
Step 6: Meeting the Needs of the EDS Child
In making plans to meet the needs of the EDS child, a 504 Plan or ESE program should be addressed. A physical therapist and occupational therapist would be beneficial in evaluating and planning for accessibility and adaptation needs in the school.
Below is a list of some accommodation that other EDS students have found to be beneficial.
Physical education/sports needs:
- Modified physical education
- Alternative health related credit in lieu of PE.
- Restrict contact sports
- Restrict weight bearing activities on arms, wrists (such as handstands, cartwheels) etc.
- Limit exposure to the sun.
- Allow storage of ice packs/gel packs
- Allow doctor note permitting use of nonprescription pain medications
- Rest periods in the middle of the day for fatigue
- Two sets of text books: one for home, one for school
- Priority seating
- Allow use of a chair instead of floor for circle time
- Adjusted chair /table height
- Use of elevators
- Allow passes for frequent bathroom breaks
- Book bag on wheels
- Help with note taking
- Extra time to get to and from class (leave before the bell)
- Stretching or walking to relieve stress
- Extended time for tests and assignments
- Extra time on timed/standardized tests
- Eliminate handwriting grade in favor of grades for content and effort
- Lockers: assign locker at eye level, allow digital lock instead of combination locks or alternative to standard locker such as a "safe" place to store belongings.
- Copies of worksheets that allow fill in blanks or underlining in lieu of rewriting existing questions/sentences.
- Chair with arms for upper body support
- Height adjustments on chairs/desks
- Pad for chair seat or back
- Pencil grips
- Pad for sitting on the floor
- Use of computers with ergonomic keyboards or Alpha Smart
Many parents and teachers have contacted the Ehlers-Danlos National Foundation with questions about EDS, ranging from basic concerns for the safety of the student, to questions about accommodations to meet the needs of the EDS child. This booklet was prepared to answer some of the questions asked most commonly by both parents and teachers. Additional information on medical aspects of this syndrome is available from the Ehlers-Danlos National Foundation, or at www.ednf.org.
Meeting the Needs of the EDS Child
It is wise to make all school staff, including lunchroom supervisors, nurses, and office, aware that a medical emergency may occur and to outline the basic plan of action. Although most children are accustomed to dealing with the problem and may give valuable advice at the time of the accident or injury, it is best to have a prepared/parent approved plan ahead of time. Always provide current updated parent contact information.
In making plans to meet the needs of the EDS child, a 504 Plan or ESE program should be addressed. A physical and occupational therapist would be beneficial in evaluating and planning for accessibility and adaptation needs in the school.
Below is a list of some accommodation that other EDS students have found to be beneficial.
Physical education/sports needs:
Modified physical education
Alternative health related credit in lieu of PE.
Restrict contact sports
Restrict weight bearing activities on arms, wrists (such as hand-stands, cartwheels) etc.
Limit exposure to the sun.
Allow storage of ice packs/gel packs
Allow doctor note permitting use of non-prescription pain medications
Rest periods in the middle of the day for fatigue
Two sets of text books: one for home, one for school
Allow use of a chair instead of floor for circle time
Adjusted chair /table height
Use of elevators
Allow passes for frequent bathroom breaks
Book bag on wheels
Help with note taking
Extra time to get to and from class (leave before the bell)
Stretching or walking to relieve stress
Extended time for tests and assignments
Extra time on timed/standardized tests
Eliminate handwriting grade in favor of grades for content and effort
Lockers: assign locker at eye level, allow digital lock instead of combination locks or alternative to standard locker such as a "safe" place to store belongings.
Copies of worksheets that allow fill in blanks or underlining in lieu of rewriting existing questions/sentences.
Chair with arms for upper body support
Height adjustments on chairs/desks
Pad for chair seat or back
Pad for sitting on the floor
Use of computers with ergonomic keyboards or Alpha Smart
School Site Interventions
You and your child have certain rights under the Section 504 of the Rehabilitation Act of 1973. (See document Appendix # 2).
All students who are considered exceptional under the Individuals with Disabilities Act (IDEA) are also considered disabled under section 504. However, all 504 disabled students are not considered exceptional under IDEA. (See Appendices #3 for explanation of differences)
In most districts, those children considered disabled under the IDEA, will have goals and objectives written in an Individual Education Plan (IEP). These often are written to include accommodations. A 504 plan is different than an IEP, as it addresses the accommodations needed instead of long and short term goals.
Each school district handles these plans a little differently as far as referral method or request for services. Generally, it is the Guidance department who manages these cases.
Knowing your rights is important, but having a good working relationship with school personnel and teachers is often more important. Go into all meeting with an attitude of cooperation. Being combative, intimidating, or demanding is not productive. Litigation takes a long time and your child's needs are immediate. The ideal situation is one where parents/teachers/schools are working together, not pulling apart. Solving your child's individual problems may require a new or unique approach. If the plan, as written, does not meet your child's needs, or if the needs change over time, be aware that you may request additional meetings to change the 504 plan or the IEP.
You are your child's advocate, which means you are their voice. Ask you child what he/she needs most. Include the child in the plans. During meetings and discussions about your child's needs, keep your child as the center of discussion. When or if the topic sways to other focuses, draw the conversation back to the child and meeting his/her needs. Stick to the facts as you discuss EDS. Discuss options, and be open to new ideas for meeting your child's needs. When discussion is finished, ask for action.
Below are helpful suggestions from parents who have already been where you are going. Be aware that asking the school for special considerations does not mean they must provide all the items you are requesting. For example, if the school allows your child to use a rolling backpack or to store therapeutic ice packs in the refrigerator, you will most likely be providing the items.
· Check with your school system to see what programs are already in place at your school or in your system.
· Ask your doctor to provide a letter of diagnosis and suggestions for the school.
· Ask about 504 plans at your school, as well as occupational and physical therapy and evaluations.
· Make appointments with the school guidance counselor, school nurse and the individual teachers BEFORE problems arise. Make a plan for emergencies to be posted with the nurse/teacher.
· Try to educate the school ahead of time about EDS. Provide information on EDS for inclusion in the Cumulative Record. Provide an EDNF CD-ROM.
· Do not assume that teachers know about your child's needs. You should be sure they are aware of the information each year.
· If you or your child are comfortable presenting information, ask to do so for the teachers/faculty or students.
· Make a list of the problems your child is having. What are some of the solutions you believe might help?
· Be clear about the things your child should not do. Things he/she should do.
· Be a creative problem solver. Think of ways to "make it happen".
· Bring a list of possible suggestions with you to your meetings. Prioritize those that are important "all of the time", and those that "may" be needed some of the time. This list and plan can be revised at your request anytime during the year or as needs change.
· Be prepared to discuss:
- Accident/emergency plan
- Physical Therapy needs
- Assistive devises / braces/ wheelchair needs
- Assistive technology needs
- Occupational Therapy needs
- Speech/Language needs
- Class modifications and interventions/accommodations needs
- Physical needs in room/hallways
- Assignments/class work needs
- Organization needs
- Special considerations
Appendix # 1. Definitions and Terms
The following are definitions or words often used by people who work with exceptional children. The definitions are meant to help or guide you and are simplified for easier understanding and use. Different school districts and educators may use these words in somewhat different ways. You should feel free to ask for definitions of words and abbreviations being used when discussing or describing your child. Different states have specific Boards of Education rules and requirements for eligibility for services and programs.
Academic - Having to do with subjects such as reading, writing, math, social studies, and science.
Accommodation - Learning to do things differently from other students because of a handicap, impairment, or disability.
Assessment - A way of collecting information about a student's special learning needs, strengths, and interests. This could include observing the student, looking at records, or evaluations and tests.
Disability - A problem or condition which makes it hard for a student to learn or do things in the same ways as most other students. A disability may be short-term or permanent.
Exceptional Student - A student who has special learning needs as described in state and local school board rules. This includes students with handicaps, a disability, or impairment, as well as those who are gifted.
Exceptionality - A special learning need. Exceptionalities include handicaps, disabilities, or impairments. Gifted is also included as an exceptionality.
Free appropriate Public Education - The words used in the federal law, the Education of the Handicapped Act, to describe an exceptional student's right to a special education which will meet his individual special learning needs, at no cost to his parents.
Handicap - A problem or condition which makes it hard for a student to learn or do things in the same ways as most of the students. A handicap may be short-term or permanent.
Homebound or Hospitalized - A kind of exceptional student education program for a student who must stay at home or in a hospital for a period of time because of a severe illness, injury, or health problem. In order to be eligible for such a program in most districts, a child must meet certain listed requirements to quality.
Impairment - A problem or condition which makes it hard for a student to learn or do things in the same ways as most other students. An impairment may be short-term or permanent.
Individual Education Plan ( IEP) - A written plan which describes an exceptional student's special individual learning needs and the exceptional education programs and services which will be given the student under IDEAS.
Least Restrictive Environment - Part of the federal law and state laws that deals with determining a handicapped child's placement. This includes that, to the maximum extent appropriate, handicapped children are educated with children who are not handicapped, an that the removal of the child from the regular school environment occurs only when the handicap is such that the child cannot be satisfactorily educated in regular classes with the use of aids and services. In choosing a child's placement in the least restrictive environment, possible harmful effects of the child and the quality of services he needs are considered.
Motor - Use of large and small muscles to move different parts of the body. Examples of motor skills are walking, holding and moving a pencil, or opening a door.
Occupational Therapy (OT) - Treatment for an exceptional student which helps him to develop mental or physical well-being in areas of daily living such as self-care and prevocational skills, etc. This treatment is given by a licensed occupational therapist. In order to be eligible for "occupational therapy" programs and services, a student usually must meet requirements.
Physically Impaired - A kind of disability or exceptionality. The physically impaired student is one who has a severe illness, condition, or disability which makes it hard for him to learn in the same ways as other students his age. In order to be eligible for "physically impaired" programs and services, a student must meet requirements.
Physical Therapy (PT) - Treatment for an exceptional student which helps to maintain or improve his use of bones, joints, muscles, and nerves. This treatment is given by a licensed physical therapist. In order to be eligible for "physical therapy" programs and services, a student usually must meet requirements.
Sensory - Having to do with the use of the senses of hearing, seeing, touching (feeling) smelling or tasting as a part of learning. An example of a sensory skill is being able to see the differences between letters of the alphabet.
Speech-Language - Having to do with a student's ability to speak (talk), write, listen, or read. This includes understanding others and making himself understood. An example of a speech-language skill is being able to put words together into a good sentence.
Speech-Language Impaired - A kind of handicap or exceptionality. The speech or language impaired student is one who has problems talking so that he can be understood, sharing ideas, expressing needs, or understanding what others are saying. IN order to be eligible for "speech language" programs and services, a student usually must meet requirements.
Appendix #2. Section 504 of the Rehabilitation Act of 1973
Section 504 of the Rehabilitation Act of 1973 is civil rights legislation that protects the civil and constitutional rights of persons with disabilities. The law states that "No otherwise qualified disabled individual in the United States… shall, solely by reason of his disability, be excluded from participation in, be denied the benefits of, or be subjected to discrimination under any program or activity receiving Federal financial assistance." Students determined qualified under Section 504 cannot be discriminated against based on their disability.
Students are considered disabled under Section 504 if they: (1) have a physical or mental impairment tat substantially limits one or more major life activities (learning or schooling is considered a major life activity), (2) have a record of such an impairment, or (3) are regarded as having such an impairment. All students who are considered exceptional under the Individuals with Disabilities Education Act (IDEA) are also considered disabled under Section 504. However, all 504 disabled students are not considered exceptional under IDEA.
Some examples of types of discrimination that 504 prohibits are:
1. Denial of opportunity to participate in or benefit from a service, educational program, or activity which is afforded to students who are not disabled.
2. Provisions of opportunity to participate in or to benefit from service, educational program, or activity which is not equal to that afforded to others.
3. Provision of aids, benefits or services that are not as effective as those provided to others.
4. Provision of different or separate benefits or services unless such action is necessary to be effective.
5. Selecting a site or location which effectively excludes persons with disabilities or subjects them to discrimination.
Appendix #3. Differences Between IDEA Disabilities and 504 Handicaps
Most students are educated in regular or General Education. This is where students achieve through a general education program in the regular classroom.
504 Handicaps allow for education and achievement through a general education instructional program with modifications recorded in a 504 plan. IDEA Disabilities allow for education and achievement through Exceptional Student Education (ESE) instructional programs as documented in an Individual Education Plan (IEP). Below are some of the differences:
What is it?
IDEA - Individuals with Disabilities Act previously called Education for Handicapped Children Act or 94-142
504 - Section 504 of the Rehabilitation Act of 1973.
IDEA - Disabilities (eligible)
504 - Handicaps (qualifies)
Who is protected?
IDEA - Students who meet qualifying conditions for 13 categories.
504 - Students who meet the definition of a physical, mental impairment which substantially limit a major life activity (including learning).
IDEA - Disabling condition results in a need for Exceptional Student Education.
504 - Handicapping condition requires an education as effective as that provided other non handicapped students.
Duty to Provide A Free Appropriate Education.
Both require the provision of a free appropriate education to students covered under them.
IDEA - Requires the district to develop IEPs. "Appropriate education" means a program designed to provide "educational benefit".
504 - Requires a district to develop a 504 Plan. "Appropriate" means an education comparable to the education provided to non handicapped students, requiring that reasonable modifications be made.
IDEA - If a student is eligible under IDEA; the district receives additional funding (FTE) to provide special services.
504 - If a student qualifies under 504, the district receives no additional funds.
504 provides protection from discrimination, not special education services.
The regulations are very similar for IDEA and 504.
IDEA - Consent is required before an initial evaluation is conducted.
Provides for independent evaluations.
Reevaluation must be conducted every so many years (usually 3).
504 - Only notice, not consent, is required.
Independent evaluations are not required.
Requires periodic reevaluations.
Both require notice to the parent or guardian with respect to identification, evaluation, and/or placement.
IDEA - Compliance is monitored by the State's Bureau of Education for Exceptional Students Division for Public School in that particular State Department of Education.
Complaints resolved by that same Bureau.
504 - Enforced by the US Office of Civil Rights
Complaints resolved by that same office.
Types of Disabilities and Handicaps
Listed below are examples of some usual types of eligible IDEA disabilities and 504 qualified handicaps:
IDEA - mental retardation, hearing impairments including deafness, speech or language impairments, visual impairments including blindness, serious emotional disturbances, orthopedic impairments, autism, traumatic brain injury, other health impairments, specific learning disabilities
504 - asthma, allergies, Attention Deficit Disorder (ADD/ADHD), behavioral difficulties, cancer, diabetes, drug addiction, epilepsy, heart disease, hemophilia, HIV, Sickle-Cell Anemia, Tuberculosis, other diseases/disorders and physical handicaps.
Appendix # 4. RIGHTS OF STUDENTS WITH DISABILITIES (IDEA & 504 PLAN) TO ASSISTIVE TECHNOLOGY
The following is taken form the Advocacy Center for Persons with Disabilities fact sheet 11/97. This fact sheet is not intended as a substitute for legal advice.
The 1997 Amendments to the IDEA (Individuals with Disabilities Education Act) require that the need for assistive technology be considered at the IEP meeting. By becoming familiar with the right to assistive technology, parents and student will be better prepared to advocate for needed assistive technology in the IEP, and also in the 504 Plan (equal access to all school programs), thereby promoting enhanced learning and functioning in inclusive environments.
What is Assistive Technology?
Assistive technology (AT) includes devices and services as well as training that help an individual to select and utilize a device or aid. AT devices are items, pieces of equipment or system (both off -the-shelf and customized) used to increase, maintain or improve the functional capabilities of students with disabilities.
Assistive technology services include evaluation, maintenance or repair and training for students, professionals or families. AT devices or aids include, but are not
limited to the following:
· Augmentative communication devices, including talking computers
· Assistive listening devices, including hearing aids, personal hearing aids, personal FM units, closed-caption TVs and teletype machines(TDDs)
· Specially adapted learning games, toys and recreation equipment
· Computer-assisted instruction, drawing software
· Electronic tools (scanners with speech synthesizers, tape recorders, word processors)
· Curriculum and textbook adaptations (e.g. audio format, large print format, Braille)
· Copies of overheads, transparencies and notes
· Adaptation of the learning environment, such a special desks, modified learning stations, computer touch screens or different computer keyboards
· Adaptive mobility devices for driver's education
· Orthotics such as hand braces to facilitate writing skills |
An interesting article about ancient Greek homosexuality. Its interesting- I can't vouch for its accuracy as this is a subject on which I'm woefully ignorant but have always been partly intrigued by given the many references in Plato to the practise. It turns out as you would expect that homosexuality in Greece evolved over the years- particularly by the end of the Athenian democracy you had people who were as the author suggests what we would recognise as homosexuals. Homosexuality for some men was a stage in development between asexual youth and the marriage bed- but others seemed to delight more in the company of men than of women. Its an interesting subject and if anyone knows more please enlighten me as to whether the assessment in the article is right or not.
November 10, 2007
November 09, 2007
The first Emperor of China is a historical character and his legacy defines in many ways what China is today. He originally was not Emperor of China, but the Prince of a powerful western Kingdom Qin. During his reign as King of Qin, he conquered the other kingdoms which constituted ancient China. The King of Qin became an emperor in 221BC over a vast landmass, stretching perhaps over a third of what is modern China today. His power was extensive- Chinese histories credit him with an almost totalitarian ideology, an aim of unification which stretched to the elimination of any possible rival, including the massacre of 460 scholars and the destruction of older feudal patterns of service and government. He brought in a single currency and connected together the walls that previous Chinese governments had constructed to the north, to build the first defensive Great Wall. The Emperor's dynasty lasted a very short time- within years of his death in 210BC, his son the second Emperor was killed and chaos descended before the rise of the Han Emperors beggining in 202BC.
The Emperor though left much behind him. The Han reigned to some extent in conformity with his principles especially of unity- and the shape of the currency that he had originally drafted remained the same right up until the early 20th Century. Much of our account of his acheivement comes from the Han historian, Sima Qian, who was born in 145BC and whose histories cover the whole of Chinese history from its mythical origins to his own lifetime. Sima Qian was hostile to the Qin Emperor partly because his dynasty replaced that of the Qin, and his history is not a history as we would recognise it in modern terms. Sima Qian writes fables and chronicles and treatises on subjects, the past for him is a set of exempla and a set of dates. He doesn't dwell as we might like him to on subjects relevant to us, but rather has the preoccupations of a Han civil servant: so his book tells us of stories about assassins, stories about how to govern and how not to govern, chronicles of dates and all from a perspective that denegrates the Qin. Despite that Sima Qian is one of the great historians of the ancient world- his name deserves to be up there with the great classical historians.
However we are incredibly lucky when it comes to the Qin Emperor, for in the mid-1970s a peasant in China came across a stupendous find. In the soil his spade hit a terracotta head, and archaeologists coming across to work on the site found not one but thousands of terracotta bodies and artefacts scattered in the soil. Having reconstructed what the site must have been, they worked out that these terracotta bodies constituted a seperate state that the Qin Emperor hoped to rule in his afterlife. At the British Museum in London at the moment, some of those finds are being exhibited. You see all sorts of people that the Emperor required in his afterlife: he has strong men, acrobats, musicians, civil servants, soldiers of all types and even a royal charioteer. Some of these artefacts bring to life stories from Sima Qian's accounts. For example on the Emperor's death, his senior civil servant Li Si kept the Imperial demise secret. He did so by maintaining the illusion that the Emperor was still alive giving orders from his Imperial chariot- and to some extent when one sees the chariot, one can imagine how that worked. The Emperor closeted and secretive and Li Si and a couple of others conspicuously running in and out to receive orders.
The terracotta army itself is shown in all its glory. It is incredible what the craftsmen (probably conscripted) could do. The skill with which the faces in particular are rendered is stunning- the visual impressiveness of what you see makes you reel back, considering that these are faces looking straight at you from thousands of years ago. The picture in particular of a fiery Turkish looking light infantryman stayed in my mind all of last night. The Museum have organised the exhibition in a very proffessional way- first they show you some Qin artefacts and describe the role of the Qin Emperor in Chinese history, avoiding much of the detail but trying to give a non-sinologist a good understanding of what this man was and what he represents. Then you proceed to see the terracotta army and court itself- which is a stunning experience and having it put in context before you see it, it becomes more impressive. The Emperor constructed this army to protect him in his afterlife- it appears they were stationed on the only open access route to his tomb in order to guard it. His tomb itself has never been opened and apart from Sima Qian's fantastic descriptions and some scientific work above the site on concentrations of metals found underneath, noone knows what is there. What we have though is these soldiers- we know they were painted and so their rather mundane colours today aren't as impressive as the gaudy way they were decorated- we know that irises for instance were painted in the eyes and we can tell all this thanks to chemical analysis of the surface of the statues. They are beautifully vibrant and vital. Each has its own character and facial expression, beard and overall look.
China is one of the hardest societies I have ever tried to understand. I have only been there once- but that's once more than most Westerners. Reading its literature and looking at its art is a very foreign experience in the way that reading Islamic literature or even Indian literature is not. Through accidents of history, China seems like another region of the earth from Europe. But its an increasingly powerful and important place- from films by great directors like Zhang Yimou to its economic importance, China is not merely an object of curious interest for the West, it is a place we have to understand. This exhibition therefore is a wonderful opportunity to learn something about China and the way that it was created and its history. The terracotta warriors are so impressive that they are a reminder of the grandeur of Chinese civilisation. They are also an incentive because of their beauty to try and understand more about the culture from which they sprang, seeing their beauty inspired me to buy translated fragments from Sima Qian's history. An exhibition like this is precisely the thing that the world's museums should increasingly engage in- if there is to be dialogue between our cultures then this is a wonderful way of expressing it and I hope some British treasures make their way temporarily to Beijing.
The Museum's exhibition reminds one of the importance of Chinese civilisation and the importance of cultivating an understanding of it. It also reminded me very visibly of the difficulties of historical research. There is so much that we do not know and will never know about the first Emperor. The history that we have is fragmented and written long after the Emperor's death. We have these artefacts but with many of them we are not sure of their use- and we have not yet seen inside the tomb of the Emperor to see what clues lie there.
One thing I do regret about the museum's exhibition is that there was not more outside or inside from historians of the era, Chinese and Western, discussing the Emperor. There wasn't even a good academic biography for sale- an unpardonable lapse! Another gap was that the First Emperor's attitude to religion was left untouched. We were invited to see the army as a simplistic guard for the afterlife or as a manifestation of the Emperor's meglamania: but I would have liked to see something more about what Chinese people of that time beleived about the afterlife and how that connected to what the Emperor did. One interesting question that wasn't touched upon was why none of his successors made this kind of tomb- it could be that they did and the tombs are lost waiting a farmer to discover them, it could be that his example discredited the practice, it could be that beliefs had shifted, it could be that this is one of many such tombs, leaving the exhibition I was none the wiser. One felt like screaming for more information. But having said that, that is possibly the churlish attitude to take. The exhibition is wonderful- the fact that these statues have left China must have been a great diplomatic acheivement and the museum has arranged them suitably well.
The First Emperor is one of those figures whose actions had momentous consequences spreading out through time, doubling and redoubling until his creation, a unified China, became one of the great powers of a globalised world in the 20th Century. Seeing the terracotta warriors, seeing the artefacts he collected around himself in his afterlife, one gets a sense of the immense power that he wielded, the creative wills that bent to his commanding will and the strength of his shortlived imperium.
November 08, 2007
to race cars. Lord Drayson has fallen on his sword in order to join the Le Mans race. I have to say that I have no idea about Drayson's record as a minister but as soon as I saw this, I rejoiced, long live the politicians for whom the hinterland matters more than the greasy pole!
November 07, 2007
Mark Steyn has a way of shocking me by producing some really good articles at times- I think he does this out of spite, he knows that I don't like some of his work and he wants me to be spinning in confusion unsure whether to like or dislike him. Sorry my sense of humour got the better of me tonight!
Anyway today Steyn has produced I think an excellent article about popular music and the need for a canon. It is really a wonderful defence of learning for the sake of appreciation. Basically Steyn's point is that you can't understand why the Beatles are great unless you understand why Bach is great. The two go together- to understand the one is to understand the other. He makes a point about the way that in order to understand something's greatness, you have to be able to see it in its context, to see what developed around it, why that move was important. Its crucial that Picasso could paint landscapes and had been trained because then his other paintings developed a meaning, its vital that Duke Ellington could play the classic solos because then he could use them in his own work. I agree completely with him: one of the wonders of artistic knowledge is the way that it supports itself. Every time I watch a new film, or read a new book (those being the two art forms I know) they tell me something about all those previous artworks I've seen and watched. And there is a strict heirarchy of knowledge in art- I would listen to Martin Scorsese for hours on film if I could because he has watched everything, and has interesting ideas about all of what he has seen.
Music is something sadly on which I'm not able to comment. One of the most illuminating moments of my life was sitting with a friend who understood music in a jazz bar in Prague. He described to me the way that what I saw as a cool sound, was actually the product of a complex interweaving of notes, a lattice of harmonies. Suddenly I saw music for a moment as this beautiful structure, which people played with, understood and manipulated- suddenly it became more than a simple nice tune, it became art, something I cared for and might grow to love. I think that appreciation is to be valued. It isn't easy to get to- appreciation of the arts is a real cost. Its something that takes time and effort, its something that you have to struggle to get to and it is something that relies on context. To take writing, its because I understand the history of English poetry that I can appreciate the opening line of the Wasteland, that April is the cruellest month- in that opening line Elliot tells us that everything that has gone before resting on Chaucer is wrong. That April is not the month of gentle showers but the month of cruelty. Poetry and novels are echoing always with previous works- the anxiety of influence was a disease that Harold Bloom diagnosed flowing through each and every author.
The great writers though manage to combine that with accessibility. I learnt to read novels- and I have to say watch films (the great twentieth century entertainment) because I began through enjoying them, I ended appreciating the same books. Most of the early readers started the same way, Jonathan Rose writes illuminatingly about the way that the first Labour MPs for example read Dickens, Carlyle, Ruskin and others and thought about them in their own way. There is a wonderful novel which really describes this process which unfortunately I can't lay my hands on right now- as soon as I find my copy I'll review it- but what shines through that book is the importance of embibing cultural classics to discovering the world of culture. The route to Austen is the route through Austen, the same goes for all the great writers and indeed for filmmakers from Orson Welles and Michael Curtiz to David Cronenberg. Its when you are bitten with the bug that you know that you have fallen in love and through falling in love you learn to appreciate and to link everything together and understand this lattice of things which all have been created partly for your pleasure.
Steyn is entirely right- you can enjoy the arts (I enjoy Music in this sense) without knowing much, but you enjoy them a hell of a lot more when you have exposed yourself to even more. Part of life is a continual adventure in self improvement- I definitely think that there are 'miles to go before I sleep' and probably will be when I'm dead- and I think that goes for art as well as anything else. There is always something 'further up and further in' to look at, there is always something which can prompt you to understand more or to reevaluate what was once familiar and now is strange. Sometimes I think in modern life we are too comfortable, the truth is that life is an adventure of understanding. For us who lag, it is worth looking up to those who are scaling the heights, but if they are worth looking up to then they are looking in admiration at the next climber. Nobody arrives at the summit, but the effort is what makes everything worth while- because by mastering that interesting novel you suddenly have another angle on human experience. Sitting down and saying no further is surrendering that knowledge and beauty that you might acquire by going up another notch- the world is limitless and its beauties are vast.
Steyn is right. To step back is folly, to stop is folly, and in this quest the canon (the works judged before by others as good) is a useful if not flawless guide. Relaxing in a comfort zone of the works written in your own culture or your own time is a waste- there is more to see and life is too short not to read that Egyptian novelist, see that Iranian film, find out about that twelfth century monk's poetry and listen to some Beethoven before going to watch Belle and Sebastien.
November 06, 2007
At the end of the First World War, the great empires of Eastern Europe, the Russian, Austrian, Prussian and Ottoman all collapsed and were replaced with a variety of successor states. Some of those states were carved out by the treaties like Lausanne and Versailles after the war, others were essentially created by military facts on the ground- and in most cases the treaty recognised what had already happened. Its worth remembering that most of the territorial changes in Europe occurred far away from the areas in which the dominant powers at Versailles- the US, UK and France- had their troops- ie the North East corner of France. Look at a map of Western Europe in 1914 and the frontiers haven't changed really that much up to today, look at a map of Eastern Europe and the world is completely different.
What happened in 1918 in order to accomplish that, and happened in 1945 as well, was the massive transfer of populations across frontiers. We often think of that as a fairly harmless process- it wasn't. To take one example, for centuries, for millennia, numerous Greeks had lived in Asia Minor. Thales one of the first philosophers, if not the first, lived for example in Miletus on the coast of modern day Turkey. By the time of the Ottoman Empire, those people calling themselves Greeks still lived there- still constituted a large minority in cities like Istanbul, Smyrna and other places. In the period after World War One the Greeks and Turks battled over the frontier between their states, in 1922 the Greeks finally lost and withdrew from Asia Minor and as they did, the Greeks living there were forced out as well. I thought of this when I first heard of it, doing my history GCSE, as a fact of history, a bloodless fact- in fact of course it wasn't- there was great brutality.
Just to appreciate how horrible that process of ethnic movement was, its worth looking at some of the accounts from Greeks at the time. Thalia Pandiri has collected some and published translations in the International Literary Quarterly- I suggest you go and have a read, but what she describes is truly horrifying. Women with sticks driven through their bodies till they emerge coming out of their mouths. Some of the stories are equally horrifying for the poverty they display- women feeding children flour in water for example or walking for miles with a bag gripped between their teeth and a child in each hand. When they arrived in Greece, many of them found a less than hospitable reception awaiting them as well. Many of them afterall looked not to the new Greece but to the Russian Tsar, traditional protector of orthodox Christians in the Ottoman Empire, as their prince.
Bringing up old atrocities has more purpose to just wallowing in misfortune. The experience of Greeks moving from Asia Minor to European Greece was horrific, but it is relatively unknown. It highlights something though of worth to consider- that moving populations is always difficult. You encounter the fact that people don't want to leave their homes, you encounter the fact that newcomers aren't always welcome when they arrive. That is even true, when unlike say in Palestine, the moving population are in the end absorbed by another population- as in the Greek case where most of the immigrants report that they did eventually become successful Greeks. Ultimately though the experience of the Greeks moving across from Asia to Europe reminds us of two things: firstly that we should not be blase about moving populations around the globe- should for example climate change result in the destruction of Bangladesh we would see the events of Asia Minor on an even greater scale even if we found somewhere for those people to go. Secondly and perhaps more importantly, it reminds us of our own powerlessness. By the end of World War One, there was barely an army around apart from those of the Western Allies and even then in Eastern Europe, it was the facts on the ground that mattered, not the pious declarations from Paris, London and Washington. International politics requires modesty as well as ambition.
November 05, 2007
An interesting piece in the Guardian reports comments from the MI5 head, Jonathan Evans, that increasingly Al Qaeda is targetting its recruitment efforts at younger and younger Muslims. In particular the organisation is looking to young British Muslims in their teens. Obviously the teenage years are amongst the prime years for people to form adult identities. One of the issues surrounding that is that people in their teenage years are often uncomfortable or unsure about where they are and what they are. They are thus prime for recruitment by groups like Al Qaeda which offer a strong identity and a purpose to life at a time when most people are going through confused emotional tempests.
Part of the problem of course is what we do about this- ultimately it comes down in part to a working education system which isn't segregated (segregation is a wonderful way to manufacture resentment from afar). No doubt, youth workers, youth organisations, parents and mosques (as well as a host of others that I've forgotten) can help as well but the spectacle of the teenage suicide bomber may grow depressingly familiar as we go into the future.
Well the Liberal Left Conspiracy came to the internet today- obviously as a group it has existed for a long time- liberals and lefties, the gay mafia, the illuminati and the free masons not to mention commies and various others have been conspiring for years which is why they have been quite so successful on both sides of the Atlantic in maintaining their control over the world. I am one of the conspirators as anyone looking at the roster will know- and I have to say I'm proud to be. The right has organised brilliantly on the internet- and Conservative Home is a really good clearing house for rightwing ideas- I know some of the best rightwing bloggers like say Matt Sinclair have written there. There isn't really any equivalent place to meet leftwing people and discuss politics on the net- Labour home is not as good as Conservative Home, its often too insular and focused in on Labour party internal affairs, other places are dominated by different sectional interests- its time the left came together in the UK on the net- and this is one option, lets hope it succeeds for doing that.
Ok lets turn to the whole idea of the liberal left- what does it mean to be on the liberal left and why do those words fit together. Lets define them first: broadly speaking I think that to be on the left is to be concerned about equality, and that to be liberal is to be concerned about freedom. The point about equality is that it produces freedom. Wealth is power- money would be nothing unless it had a value and that value is the goods and services it commands. The more wealth that someone has and the more independent that wealth from the interference of others, the freer they are to gain what they want in life. Rightwingers believe that the only obstacle to a free will is a state: they are right that the state can be a significant obstacle to the exercise of a free will, noone with any knowledge of this century could deny that and many on the left stood against the state as it limited the freedom of will (Orwell is a great example) but rightwingers are wrong to say that it is only the state which obstructs freedom. Corporations do too- and even the wealthy can obstruct liberty- both can use the state as well in their own interests- you could argue that that is what the British libel laws do.
Equality is married to freedom thus at a fundamental level- because without equality I cannot be free. Its encapsulated in that old piece of wisdom that beggars can't be choosers- something that the right tend to forget. This isn't an argument for state socialism, it could be but it isn't. It isn't an argument for any particular vision of society. But it is an argument that you cannot have real freedom without having equality, that you cannot be concerned about liberal things, without being concerned about leftwing things. And that goes as well for many of the other battles that the left are involved in, freeing women from the dominion of their husbands, freeing homosexual people from the legal restrictions of those that don't share their morality, freeing the innocent from the tyranny of a despot who would rather hold us all in jail than listen to any of us. All these things are both leftwing and liberal- how they are achieved is a totally separate issue but they can only be acheived if we think about equality and freedom together and try to acheive both through our policies.
That's why I'm conspiring for the liberal left (though I have to say this blog will remain basically what it has always been)!
November 04, 2007
Clive James is a figure unlike most others in our world- James has made a career of being an omnivore. From the chatshow couch to the comic circuit to the learned essay, James has succeeded everywhere he has gone. Writing and broadcasting, he has turned his natural wit to good account and provided a series of sparkling memoirs to furnish the bookshelves of the learned with. Cultural Amnesia, his latest book, is a fine effort to capture the unique folds of James's own mental landscape- he provides a short essay on over 100 cultural characters mainly from the last century. All the essays come out of a single quote- and often James doesn't even pause to ponder the life, instead pondering the importance of that quote.
The quoted range from Duke Ellington to Hegel, Federico Fellini to Margerate Thatcher, from Tacitus and Edward Gibbon to Coco Chanel and Adolf Hitler. The range is astonishing- though the absense of any scientists is equally astonishing. James mentions an Albert Einstein but its the musician not his more famous namesake and relative the physicist. Indeed science is one of the leading absenses from the collection which is biassed very much towards the arts. Analytical philosophy is also underepresented- we have an essay on Wittgenstein but characteristically in it philosophy students are dismissed for giving him the 'credit for everything that would have struck them if they had ever been left along with the merest metaphysical lyric from the early seventeenth century.' The Wittgenstein that matters to philosophers is the one that 'they can prove only to each other' and what James is interested in is the Wittgenstein that matters to the writer- to the humanist.
For that is what this book really is, a monument to what we might call humanism. A humanism that sees the limits of the human as surely as it does the extent of his range. James is limited- but to stress that is to undermine really his acheivement here- which is to gather and express particles of knowledge and understanding across many fields and many languages. He gets some judgements wrong- he dismisses Edward Gibbon as a poor stylist. James tells us that 'what he [Gibbon] wrote rarely lets you forget that it has been written'- possibly that's true but its also Gibbon's virtue and not to see that is to miss what Gibbon was trying to do and therefore to criticise him by a standerd he wasn't attempting to reach. James doesn't get Gibbon's historical breadth or depth either- doesn't see that the styllistic tics are made up for by the fact that Gibbon was another such as James who spanned centuries in a massive project that will probably never be attempted let alone completed again.Quotation has this feature that it inspires you to seek out the epigram- the fragment that illuminates rather than the rolling cadence of prose. Martial the great Latin poet is perhaps the most eminently quotable of Latin poets in that what he wrote was bitchy and short, James in these essays has the same quality. Like the greatest essayists he can skewer wonderfully. He can also at his best capture real nuance- his description of Edward Said in this sentence is perfect, 'As a critic and man of letters he has an enviable scope but it is continually invaded by his political strictness'. It captures the many sidedness of Said- the political lack of nuance which led him to some cartoonish descriptions of orientalists and of the orient but also the greatness- for Said who always recognised Israel and wanted Palestinians to recognise the sorrows of the Jews was a great man. James is able to capture that and through a quotation of Said's about the Battle of Algiers, bring to life the double sidedness of Said.
But this book is not all nuance. James is more often than not on the good side and vows war against those who cravenly boosted tyranny. He writes eloquently about the Manns- Heinrich, Thomas and Golo- all of whom resisted Hitler from outside the boundaries of exile. Of all the praise though it is that devoted to Sophie Scholl which most resonated with me. Scholl, James tells us, 'was probably a saint' and died in complete silence. What James wants to do with praise is make us think- he points to the fact that in his judgement despite the fact that there is a perfect actress for the role alive today (Natalie Portman) Scholl should never be portrayed by Hollywood. The finality of her end is her tragedy- far better for it to be a more obscure German film starring the unknown Julia Jenstch to portray her for the public so that they too understand the finality of the fall of the ax upon her neck shut out one of the true heroines of the twentieth century and sent her to darkness.
If Scholl volunteered to die, despite the fact she did not have to, to make a point against an odious regime, then James rightly eviscerates those who have supported those odious regimes. Though Sartre is his betenoir- he hates Sartre's evading of responsibility, hates the fact that 'Sartre was called profound because it sounded if he was either that or nothing' but ultimately his essay on Sartre is not the most interesting. Rather I think it is the essay on a much slighter figure- Peirre Drieu La Rochelle- a leading intellectual of Vichy that really made me think. For what he captures in that essay is the moment of victory in 1945, when the Germans were driven out and La Rochelle committed suicide. The key fact for James though is to evaluate the hysteria- a hysteria he informs us drily that Sartre backed and that Camus (who actually had a resistance record) disdained (though Camus thought there ought to be a reckoning). He leaves us in no doubt of the guilt of Pierre Drieu La Rochelle- but also paints a picture of France in those years which is terrifyingly accurate.
Totalitarianism is one of the foci of this book- James argues long and hard against it. Whether it is Communist or Fascist, he suggests it is deeply repugnant and you get the sense that he thinks that clear writing, thinking and reading are its enemies. As he said recently to Stephen Colbert, intellectuals get things wrong all the time- but they get them wrong less than those who don't open themselves to intellectual pursuits. In reality this book is a book about heroes- but it is not a book about heroism. The essay structure enables there to be a convincing absense of structure- in the sense that James is not interested in archetypes but in individuals- his essays are at their most effective when they describe either of two things- the impact of writing upon him as an individual or the way that this individual's career worked. An essay on Nadezhda Mandelstam is incredibly effective at making you realise the pain that she must have felt as the Stalinist machinery of death whirled past her windows. It drives you to the reality of the statistics.
Though James is reassuringly committed to the dry substance of the real world, he is most acute when he focuses on individual experiences, exploring them and rendering them to his reader. His selection is driven, as he argues in his essay on Chris Marker, by the solidity of the facts that he sees and understands but his talent is for explaining experience. This is a book which is unashamedly focused on reality- James gives postmodernism and its creeds of unreality very short shrift indeed. He is openly contemptuous of philosophical relativism and disdain for truth- openly praises the empirical and solidly researched. He bases his love for art upon a respect for reality.
James's range of understanding in this book is incredible. James is a great evoker of what other authors do and write and film and play. He can convey the meaning of others' statements in such a way as to make you want to read and listen to and watch their books, music and films. He makes you want to stroll down the streets of Vienna in particular and pop into the cafes to hear the arguments and consume the culture. He makes you want to open the books, to understand what Contini means when he says that you need to learn poetry. He creates a desire in you to leap from cultural tree to tree- as James himself in these essays does- referring for instance in an essay on Marc Bloch to the seductions and disappointments of Pound's poetry. He made me want to learn languages- to read these authors in their original tongues and capture the calligraphy of sound that they all employed.
Ultimately there isn't a greater compliment for a book like this than to say that- to say that this book is like the trunk of a great tree, along whose branches if you pursue them are fruit much more gaudy than anything found in the original bark. This is a book that leads to other books. Its a book that can be read at one sitting or dipped into- yes there are mistakes and there are manifold errors. But to forgive someone for misunderstanding that Gibbon is amongst the greatest English historians requires a great acheivement and this book is a great and interesting acheivement.
The other day I wrote an article on analytical blogging, which got some negative attention from Dizzy, who makes a fairly amusing point against it though personally I'm not as convinced as he is that intelligence is only reserved for the elite. Its interesting as well that modern conservatives often have tended towards being unabashedly in favour of populism- that reinforces one of my feelings that modern conservatism and other historical forms of conservatism are not the same- I can't imagine Edmund Burke or Hayek even giving three cheers for the Sun in the way that Dizzy does!
However that isn't the main point of this post. Matt Sinclair asks a much more interesting question about smart people and blogging, and I think he is right to ask it and the answer in the case of this blog demonstrates something which I think is interesting. Matt asks "Why should someone with interesting and novel things to say use the blogosphere as a medium?", he goes on to deliver some interesting answers, all of which depend mostly on the community as a whole providing a forum. Matt imagines that blogging is a bit like an intellectual salon on the net, in which we can throw around ideas, as he rightly points out that presumes a membership, there is no point talking to onesself.
Somebody asked me on my thread about this, why I don't do more analytical work on politics. I do a bit, but nowhere near what Chris Dillow does on Stumbling and Mumbling- and I think this ties into another reason to maintain a blog, which is one of the basic reasons that Westminster Wisdom (the title is partly ironic) exists. This blog really isn't an analytical policy blog- though I do occasionally rummage through politics and policy, its really a purely egoistic exercise. For me a blog is the equivalent of an 18th Century common place book, ie its where I put down my impressions of the world so I can go back to them. An interesting quote, a fun video, a film review, even a review of a novel, anything which makes me remember how I reacted to something for the first time.
I think that is a valid reason to keep a blog- partly because experience flows past me at such a rate that I can never really grab hold of it. Throughout my life, amongst my major vices is forgetfulness, and that means that I often lose hold of what I should know or should remember. Here I have a resource to which I can turn, when I want to, to find out about say Rousseau's walks or Bresson's Joan of Arc. Part of that is it forces me to think about what I see and read more acutely than ever before: because I know I'm going to have to write an article up here on it. That makes me look deeper and try and understand more. Its also a good resource to remember what an idiot I am occasionally- there are moments on this blog where I know I've been a complete fool- reminding onesself of that is a good thing and doing it on a blog is fairly harmless. (Which in a way brings me back to Dizzy, acute mockery of your own pretensions is always a good thing to read!)
In answer to Matt's question therefore- I think there is another reason- in addition to the good ones he has given- for a person to keep a blog and that is as an online diary. Afterall that is what blogs started off being- and I wonder whether in the end that will be their principle use.
LATER Incidentally Dizzy should probably go and watch this. |
Phil/FA 310 Introduction to Aesthetics
Instructor: Dr. Norman Lillegard Office: H 229 587 7384
Office Hours: 10.00 ‑11.00 a.m. and 1-2 p.m. MWF and by appointment
Texts: Puzzles about Art (ed. Battin etc.) and Basic Issues in Aesthetics by Eaton (In UC and Bradley).
Course Title: Questions about the Arts That Won’t Go Away.
Art; can you define it? Are there non-subjective standards of criticism/evaluation? Why bother with the arts? Why spend (public) money on galleries, subsidies for playwrights or orchestras, etc.? Are there correct interpretations of art works? Does interpretation depend upon the artists intentions? Should art be “moral” in any sense? Should it be censored in any way? Is avant-garde art really art at all? What is the “ontological status” of various art works? How do art works differ from products of the crafts, if at all? How if at all do “fine arts” differ from entertainment? Does, or should, art “imitate” life or anything else? What about “art for arts sake?” Is art a medium of expression? Expression of what? The Artist’s soul? Her emotions? Should “works” produced by computers (or monkeys!) be counted as works of art? Do the arts have a future or are the fine arts done for, about to be replaced by decoration, entertainment, craft? And so forth.
The Purposes of this Course: To address the questions listed above and related questions. To become familiar with some of the principal answers that have been proposed from the Greeks to the present. To gain practice in thinking critically and with a sense of the options available about these questions. Perhaps, to formulate some defensible views on these matters that will inform one’s own practices in the future. Perhaps, to find ways of defending civilization against barbarism (!!??).
We will be studying the views of some major thinkers, but the aim is not that you be able to repeat their views, but that you learn to think with them. Therefore, the ability to parrot views (whether those of an author, the instructor or anyone else's) or regurgitate information (like a quiz show participant) is of no use to you or anyone. You will not be tested on such an ability. Exams are designed to test understanding of arguments and issues, and critical reading skills, rather than retention of information. However, you do need to be familiar with some relevant examples and illustrations, and some terminology.
! Attend class and participate, do the readings, do all written assignments, pass the exams. Two exams. (multiple choice, T/F, see sample exams on web page). First exam worth 120 pts, Final exam is comprehensive, worth 180 pts.
! Quizzes: there will be frequent (once a week or more) unannounced quizzes. Missed quizzes cannot be made up. Each quiz will be worth 6 – 12 points, and will consist of multiple choice and T/F questions. Total, ca. 130 pts.
! A short paper (no less than 1500 words) on a topic approved by the instructor must be completed and handed in by mid-term and revised in the light of criticisms by the final. NO PAPER WILL BE ACCEPTED AFTER THE MID TERM. 100 pts.
! One report on a UTM art event (exhibit, theatre, concert). 25 pts.
! Attendance. Regular attendance and informed participation in class are essential since (a)not everything covered in class is included in the texts (b) you will need help with this material, and that is what class sessions, and the instructor, are for. 40 points.
! Extra Credit: Don’t count on much. Carefully prepared reports or other class presentations (including student prepared debates under the instructor's guidance) can earn extra points. (Max. of 30 pts.). Presentation of actual (original) art works with discussion relating to the readings might be especially valuable.
Total points ca. 600. Normally %90 of total points gets you an 'A', %80 a 'B' and so forth, but significant adjustments for curve are made when necessary.
The instructor’s web page for this course will include sample exams, lists of important terms, and outlines of classes. All quizzes will also be preserved on that page for review purposes. Access the Phil/FA 310 page through the UTM page (click on faculty staff, then on faculty web pages) or by using www.utm.edu/~nlillega/lillegard.htm. and clicking on relevant link.
Class Conduct, Instructor's Role, etc. What I Expect of Students.
1.Treat each other with respect. 2.Treat the instructor with respect. 3.Do not talk unless called on.
4. Do not leave the room without permission except in extreme emergency. 5. Be on time.
6. Be eager to learn. The best indication of progress is engagement with the issues and ideas we deal with.
7. Do not be afraid to say "I don't understand." 8. Expect the same of me as I expect of you. (Except for #3, and #4, of course. You will see that I follow #7 a lot.)
NB. Any kind of cheating is a serious offense and will be dealt with accordingly. It also ought to be beneath the dignity of each and every student.
Classes will consist of a mix of lecture, discussion, possible occasional reports, and viewing of videos,
listening to recordings, etc. Students are expected to treat other students in a polite fashion, even though they should
feel free to express disagreement on ANY topic or ANY claim that is advanced by anyone, including the instructor.
At the same time, each student must attempt to exercise responsibility by keeping discussion focused on the subject at
hand and by listening carefully to the responses of the instructor and other participants.
Particular value is placed on argument, as opposed to mere expression of opinion. Say what you believe, but be prepared to say why. The instructor will attempt to clarify difficult concepts and passages in the text, and provide relevant examples. Students are encouraged to provide their own examples. If there is an art work (painting, pot, poem etc.) by yourself or someone else that you want to discuss, bring it or describe it and some class time will be spent on it, provided there is some connection to the readings.
Students should feel free to interrupt with questions or comments, even though on occasion answers may be postponed for the sake of coherence. The instructor is pledged to careful consideration of any view, including those which he finds unsupportable, and to critical thinking with any student who values thoughtful discussion. Students who feel a need for individual help should feel free to ask..
NOTE: "Any student
eligible for and requesting academic accommodations due to a disability is requested
to provide a letter of accommodation from P.A.C.E. or
COURSE OUTLINE: (Approximate. Content and time periods may vary slightly.)
Week 1 (1/18) Overview of course. Kinds of “arts.” The notion of “fine art.” Defining art. Art criticism. Read Battin, 1-26. Be familiar with the 24 cases described.
Week 2 (1/24) Week I continued. Defining “art” and theories of art. Eaton, ch. 1.
Week 3 (1/31) Week II continued. Creativity. Selections from Battin. Eaton, ch. 2.
Week 4 (2/7)Week 3 continued. Eaton ch. 3. Viewers, taste, emotion etc. Selections from Battin ch. 2
Week 5 (2/14) Eaton ch. 4. Languages of art.
Week 6 (2/21) Eaton ch. 5. Art objects. Battin selections.
Week 7 (2/28) MIDTERM EXAM Wed. Mar. 2 Paper due. Interpretation and criticism. Eaton ch. 6. Selections from Battin ch. 3,4, 6.
Week 8 (3/7) Week 7 cont. Battin ch. 6 Forgeries, etc.
Week 9 (3/14 – 18 SPRING BREAK)
Week 10 (3/21) Week 8 continued. The value of Art. Eaton ch. 7. Selections from Battin ch. 5.
Week 11 (3/28) Meaning and truth in Literature. Handouts. Art, public policy, Eaton ch. 7
Week 12 (4/4) Art, ethics. Eaton ch. 7. Handouts. .
Week 13 (4/11) The end (or future) of the arts. Handouts.
Week 14 (4/18) . Discussions of papers, themes.
Week 15 (4/25) “ “
Week 16. Classes end May 2. FINAL EXAMS, MAY 5-11.
Article for week 13, 14. http://www.eurozine.com/article/2003_02_25_lillegard_en.html
Plato – Republic, Ion.
Aristotle – Poetics 1-15.
Mo Tzu, Hsun Tzu (On music)
Plotinus – Enneads I.6
Hume – Of the Standard of Taste
Kant – Critique of Aesthetic Judgement
Friedrich Schiller, Letters 26-27
Hegel – Intro. to Aesthetics 1-3.
Schopenhauer, World as Will and Representation
Tolstoy – On Art.
Coomoraswamy- The Dance of Siva
Dewey – Art as Experience
Heidegger – The Origin of the Work of Art
Collingwood – Principles of Art
Cavell – The Claim of Reason and Must We Mean What we Say (various essays).
Danto – The Philosophical Disenfranchisment of Art.
Wollheim – Art and its Objects
Expression in __________ (music, painting etc.)
Censorship of the Arts
Fiction and truth.
Meaning and Intention in Literature
Tragedy and Ethical Criticism of Art
Application of selected aesthetic theories to __________ (Drama, Painting, Film etc.)
Art vs. Entertainment
The Aesthetics of Nature
The Historical Nature of the Arts.
Art, Propaganda, Politics.
Battin, 1-26 Questions 1 - 21 below
1. What does case 4 add to case 1 and 2 and 3, if anything?
2. What does 5 add to 1,2,3, 4, if anything?
3. What does 6 etc.
4. What does 7
QG 2. Read Eaton, 1-13
1. What are the four general types of aesthetic theory?
2. Give some reasons for thinking aesthetic theories are not even possible.
Eaton p.14-31. 1 - below
1. Assuming that something is a work of art only if someone works on it in some way, be prepared to illustrate each of the following views about what kind of work is required to produce an artwork, as opposed to simply a (non-artistic) artifact:
a. art works result from the particular unique personalities of artists
b. art works result from a creative and original activity, where what makes an activity creative is either that it
i. consists in something more than following rules to produce an envisaged end
ii. consists in paying attention to features of things which have “artistic potential.”
c. art works result from carrying out an artistic “intention” (particular kinds of intention make certain things works of art)
d. art works result from (successful) attempts to express something (like an emotion).
2. Give criticisms of each of a-d above, as found in Eaton.
3. How does Aristotle respond to the claim that the arts are “irrational?”
4. Be prepared to illustrate and critique each of the following views on art as expression of an emotion (E): in a work of art
a. the artist expresses his/her own E in the work
b. the work arouses an E in the “viewer”
c. a and b combined: the artist expresses an E and the same E is aroused in the “viewer” (Tolstoy-type view of art as “communication”)
d. the work “depicts” an E
e. the work has in itself the properties or traits of people who feel E (Langer)
f. the work “treats something in a way that demonstrates” E (Sircello)
5. Compare and contrast Croce’s and Dewey’s views on art as an expression of the artist’s “idea.”
Eaton, p. 34-52
Battin, p. 28-58
1. Eaton suggests that there are three sets of criteria that might be used to sort people who are having an “aesthetic experience.” What are they?
2. Taste on the Hume/Sibley view includes
a. special sensitivity to properties objectively present in a work
b. the idea that the perceptual faculties employed by a person with taste are different from ordinary perceptual faculties.
c. the judgements of people with taste will converge
What are some criticisms of a, b and c?
3. Are viewer emotional responses to works of art
a. real ordinary? (does the viewer feel real sadness, for example)
b. not real but special aesthetic? (Burke’s “delight”, or, metaresponse)
c. ordinary but in a non-ordinary way? (in control etc.)
Discuss each of a-c pro and con.
4. What is the problem of “negative emotions”?
5. What are some problems with the idea that the “aesthetic attitude” is essentially “non-practical” or involves “distancing?”
6. Explain and illustrate the Sibley/Hungerford view that non-aesthetic properties never entail any aesthetic properties (or, aesthetic properties are never “condition governed.” Mention the “looks/is not really” test)
7. Give some criticisms of the Sibley/Hungerford view.
8. According to Dickie, the aesthetic attitude is a myth: the jealous husband and the lighting technician are not attending in a different way, they are attending to different ____________.
Mention some possibilities.
Read (carefully!!) Eaton, 53-74
Think about these examples: Hobo signs, a musical phrase (Handel), Bach’s Nun kom der Heiden Heiland, Lover’s kiss, Jan Steen’s painting, Battin 3.13, 3.15, a poem (cf. Battin, 2.16), Dǖrer’s Melencolia I.
Ask these questions about these examples, and write down your answers:
1.Do they refer to, or contain parts that refer to, something? If so HOW is the reference achieved?
2. Do they state something? If so, HOW?
In answering 1 and 2, take account of the fact that referring and stating are usually done through language, and that actual verbal language is completely missing from most of these examples.
When you have answered 1 and 2 as FULLY as you can, note
(a)all the similarities or identities between the different cases, and
(b) all the differences.
Remember that these answers require reference to certain fundamental contrasts, such as that between conventional and the natural, between reference and resemblance, between iconic and non-iconic signs. Also keep in mind that these contrasts are not always sharp.
Study Eaton 76-101.
Eaton, 104- 123. Battin, 60-102
Eaton, 125-147 Battin, 148-178
T or F
1. Something might seem to count as a work of art because it is displayed in a certain way.
2. One question that arises in connection with “sound sculpture” is this: are there genres of art that are constrained by the “materials” used?
1. Artist centered theories of art could include expression theories.
2. Weitz denies that it is possible to define (give necc and suff conditions for)the concept ‘work of art’ because
a. it is an “open concept”
b. there are no resemblances at all between some art works and others
c. there are at best family resemblances between some art works and others
d. a and c.
1. According to Sircello, a painting could express a love of animals in the sense that the artist painted the animals “lovingly.”
2. If a work of art expresses sadness in the sense that it made the viewer feel sad, then it is mysterious why most people would want anything to do with it.
1. The idea of the “aesthetic attitude” typically includes the ideas that
(a) a special faculty is required for aesthetic perception
(b) aesthetic response is non-practical
(c) aesthetic responses are to conditioned-governed properties
(d) all of these.
2. Dickie claims that what makes a response aesthetic is what sorts of things the “viewer” responds to.
3. Twains two different responses to a sunrise on the river illustrate
(a) the difference between a practical and a non-practical response
(b) the difference between an aesthetic and a non-aesthetic response
(c) the difference between a response that shows taste and one that does not.
(d) a and b.
1. An 18th century garden, like that at Stowe, that leaves certain things out (like a statue of Queen Ann) and locates certain things lower than others, seems to make various statements, including negations.
2. Both Gombrich and Goodman emphasize the conventional aspect of artistic representation.
3. The more a sign looks like or in some way resembles what it refers to, the more “iconic” it is.
1. It is not implausible to think of Brancusi’s “lovers kiss” as referring to or representing a lover’s kiss
(a) by resembling in some slight way two lovers
(b) by “standing for” a lover’s kiss by the use of slightly iconic signs
(c) by evoking feelings relating to unity, difference, bridging, inner and outer divisions, and the like
(d) any or all of these.
2. Jane Austen’s Pride and Prejudice could be said to
(a)communicate knowledge or understanding of human life and character
(b) be true to life
(c) be a good novel because it invites disinterested contemplation (the aesthetic attitude)
(d) all of these
(e) a and b.
1. Formalists argue that
(a) only properties directly observed in a art work are relevant to aesthetics
(b) the artists intention is relevant to aesthetics
(c) such traits as shading (in a painting), thematic development (music), or structure (novel) are aesthetically relevant.
(d) all of these
(e) a and c.
2. Clive Bell and Roger Frye are probably the best known formalist critics.
1. Marxist aesthetics is contextual in a very broad sense.
2. The context of a work of art could include
(a) the historical/political situation in which it was produced
(b) the prevailing artistic traditions at the time of its production
(c) the personal idiosyncracies of the artist
(d)all of these
(e) none of these.
1. Contextualists hold the view that all that matters in a work of art are its intrinsic properties.
2. According to the institutional theory of art, an object, X, becomes an art work
(a) when the artist produces X
(b) when art institutions confer on X the status ‘art work’
(c) when enough people see that in fact the work is worthy of appreciation
(d) none of these.
1. Everyone would agree that since James says that The Turn of the Screw is just a ghost story, that therefore the best interpretation is one that treats it as a Ghost story.
2. There is a clear difference between interpreting a work of art and evaluating it.
1. If someone interprets ‘dark satanic mills’ in Blake’s Preface to Milton to refer to sooty industrial factories, they
(a) have not taken into account the author’s intentions
(b) have taken into account the historical circumstances of the poem’s composition
(c) have obviously made the poem less enjoyable
(d) all of these.
2. Giving reasons for the belief that a work of art is a good work is not like giving reasons for believing that a breakfast is nutritious, or for believing that it is not a good idea to go 80 mph in a 30 mph zone.
1. The case of Virgil’s and Mendelssohn’s death bed wishes respecting their masterpieces illustrates the conflict between the value placed on art works and other values.
2. Tolstoy argued that the only good art works were
a. works that had economic value
b. works that appealed to people with taste
c. works that had intrinsically pleasing properties
d. all of these.
e. none of these.
3. Aesthetic value (the value of the arts) could, on some views, be
a. a function of moral values
b. in competition with economic values
c. not in competition with any other values
d. all of these
e. none of these.
1. The 18th century revisions of King Lear illustrate
a. a conflict between aesthetic and economic values
b. how some people want art works to serve moral didacticism
c. how questions of authorship can never be settled
d. all of these.
2. The belief that “earthworks” should not be subject to the requirements of environmental protection suggests that social values and aesthetic values can conflict.
3. It is clear that religious subject matter in an art work will make it a work of religious art.
Class Outlines Phil/FA 310 Aesthetics
A.What sorts of things (actions etc.) count as works of art? Or, What is the “definition” of ‘art?’ Examples:
B. Should works of art express something? Like what? Examples:
C. Is there some right way to “interpret” works of art? Examples:
1. q. What is the problem in (A)?
Answ. Deciding what features of a thing make it a work of art.
2. q. What is the problem in (B)?
Answ. Same as 1?
3. Any additional feature for (C)?
1. Organizing the intuitions:
Should we focus on
I.1 The Artist I.2The Object
(intention, (Beauty, formal
Expression etc.) properties, etc.)
I.3The Viewer I.4The Setting
A. aesthetic theories tend to focus on one of these. E.g. it is necessary and sufficient for x’s being a work of art, that it be produced with an “artistic” intention.
i. perhaps there cannot be such theories. Perhaps ‘art’ or ‘aesthetic’ or ‘beauty’ are open concepts. Weitz and Wittgenstein.
ii. language games and games.
“what is common to all games?”
a. family resemblances.
2. The cognitive status of evaluative
claims about art;
A. is ‘that is a good painting’ like
‘that is a good knife’ or not?
1. (I,1)Focusing on the artist:
A. focus on the artists inner states and the EXPRESSION of those states.
B. art works result from the particular unique personalities of artists
C. art works result from a creative and original activity, where what makes an activity creative is either that it
i. consists in something more than following rules to produce an envisaged end
ii. consists in paying attention to features of things which have “artistic potential.”
D. art works result from carrying out an artistic “intention” (particular kinds of intention make certain things works of art)
E. art works result from (successful) attempts to express something (like an emotion).
2. Exploration of E. Art as expression of an emotion (E):
in a work of art
A. the artist expresses his/her own E in the work
B. the work arouses an E in the “viewer”
C. a and b combined: the artist expresses an E and the same E is aroused in the “viewer” (Tolstoy-type view of art as “communication”)
D. the work “depicts” an E
E. the work has in itself the properties or traits of people who feel E (Langer)
F. the work “treats something in a way that demonstrates” E (Sircello)
3. Croce’s and Dewey’s views on art as an expression of the artist’s “idea.”
1. Eaton on three sets of criteria that might be used to sort people who are having an “aesthetic experience.”
2. Exploration of (1,i). The idea of Taste The Hume/Sibley view of taste includes
A. special sensitivity to properties objectively present in a work
B. the idea that the perceptual faculties employed by a person with taste are different from ordinary perceptual faculties.
C. the judgments of people with taste will converge
3. (Exploration of 1.ii)Viewer emotional responses to works of art. How should we describe those responses?
A. They consist of real ordinary emotions? (does the viewer feel real sadness, for example)
B. not real but special aesthetic? (Burke’s “delight”, or, metaresponses) (This seems to fit with 1.i)
C. ordinary (like ‘sad’) but in a non-ordinary way? (in control etc.)
D. The problem of “negative emotions” Examples:
4. Exploring another approach to 1.i.
A.The “aesthetic attitude” : essentially “non-practical” or involves “distancing.”
i. examples: from Twain, Bullough, (caught in a fog).
B. Aesthetic attitude and The Sibley/Hungerford view that non-aesthetic properties never entail any aesthetic properties (or, aesthetic properties are never “condition governed.” ‘It is red and gold’ (non-aesthetic properties) never entails ‘it is beautiful’ (aesthetic property).
i. the “looks/is not really” test.
Can sensibly say ‘it looks yellow but is not’ but cannot say ‘it looks beautiful but is not.’
7. Exploring 1.iii. Rejection of aesthetic attitude theory. Dickie: the aesthetic attitude is a myth: the jealous husband and the lighting technician are not attending in a different way, they are attending to different _things___________. (he is referring to Bullough)
Mention some possibilities.
So, What things should we attend to in works of art? Which are important?
A. Beardsley – regional qualities, formal unity. Lists on p. 49
i. examples from Battin
WEEK V The Arts and (as) Language
1. Does “getting A” (A could be any sort of non-linguistic art work) require something like learning a language?
A. In a language, words (are used to) refer to things and sentences (are used to) say things. Examples:
ques. What would correspond to words (lexical items) in non-linguistic art?
i. elements that “refer”. cf. medieval iconlogicae
ii. more general cases of “elements” “standing for” something.
Ques. What would correspond to sentences? cases of something (a string of elements) “saying” something.
2. Exploration of idea of elements in art works referring to something.
Ques. How can x stand for y?
A. by representing it? More obvious examples from painting.
Less obvious ones from music, literature.
i. representing and imitating.
ii. imitating by resembling.
iii. differences between resembling and representing. x represents a horse. Does it resemble one? Two chairs resemble each other. Do they represent each other?
iv. cf. examples from Hobo language, Handel.
v. on the other hand, a stick drawing can “represent” X. Does it resemble it?
B. If representation does not work through resemblance, how does it work? Gombrich and Goodman; it works somewhat the way a regular language does. Art really is a lot like verbal language.
i. denoting by convention.
C. The ingredients of a language – signs, syntax, semantics.
a. Signs, semantics, and natural/non-natural, or conventional, relations. Cf. ‘dog’ and a picture of a dog. Road signs.
Now, how much knowledge of conventions do we need to see what the picture “represents?” How much to produce a representation?
i. ways to classify signs. Natural. Conventional.
Another way to classify signs – icon, index(indication), symbol. Or, iconic, iconlogical, iconographical. Is a road sign iconic?
i. degrees of iconicity.
B. syntax. How the signs (“lexical items”) are arranged makes a difference to “meaning.” Cf. a sentence, a musical phrase, Steen.
3. Gombrich – reference does not depend upon resemblance.
A. making and matching, schema and correction.
i. schema are learned from the present “tradition.” Like learning a vocabulary. Thus the arts are historically located.
ii. however the individual artists can “correct” or modify the schema.
iii. figures, shapes, are not chosen because they resemble x, but because they are good “substitutes” for x.
Even in painting, the “signs” are more like symbols, less like icons, then we usually recognize.
4. Goodman – goes even further than Gombrich in denying resemblance. Shapes in paintings are conventional denoters.
A. Seems obviously wrong. In painting SEEING plays a fundamental role. Not so with “reading” a text.
B. Goodman agrees that the “symbols” in painting are not “notational.” You don’t “read the painting” and then produce it, the way you read a score, or a poem, and then produce it (perhaps many times many ways).
i. in a poem, all the identifying features of it are in the printed or spoken poem itself. Likewise, for a piece of music, Goodman thinks.
ii. in a poem, etc. the features are repeatable. That is essential according to Goodman, for being language.
iii. That is NOT the case for paintings. So they are not language. But they do employ conventional signs, like language does.
C. Goodman seems wrong about music. Can’t tell just from the score how to reproduce or “repeat” it.
D. In a language familiarity in varieties of contexts, grasp of syntax, is required for grasping the reference. Examples: rose, rose. Similarly in painting, music? Historical element. Durer.
3. Problems with the art-language analogy with respect to painting, music.
No logical connectives
No conjoining of separate works to make more complex “statements”
Not clear what the lexical items in paintings, music, are.
Truth in the Arts.
1.Different ways fiction might be true.
a. Possible worlds.
b. Fiction can tell us truths about people, places. Tolstoy on families. Dickens on London etc.
2. Pictures and truth. If no language, then no truth?
a.True to life. A picture is worth a thousand words. Etc.
b. Symbols that make a statement. Polish triangles and circles.
c. Can pictures lie? Yes. In more than one way. False in relation to some particular fact.(Curtis’ photos of Indians) Or, false to “life.”
3. Music and truth. Making a statement? Lieing?
Or is music the most formal art?
Expressing and stating.
1. The ontological status of works of art: what kind of “being” do they have? E.g. a painting, a poem, a song.
a. spatial – non spatial (temporal)
b. performing – non performing
c. works that are only types – works that have tokens.
2. What is included in the work of art? Is the context part of it? Examples.
a. Pole – the history as part of the work.
1. Context is irrelevant.
“Message” is irrelevant (art as language?).
Representation is irrelevant.
Artist’s intentions are irrelevant.
Intrinsic properties (presentational properties) are all that matter.
Ques. Is expression irrelevant?
a. examples from Painting, Music, Literature, film.
Cf. a film remake. Content remains the same. “form” is new. Cf. p. 81
Are changes part of it?
b. Ques. Is it possible to view (notice) ONLY the formal properties?
Answer: Different cases: Mona Lisa, Notre Dame, Braque, Brancusi, Liden, Strandqvist. A Bach prelude. A Bach Fugue. A Mozart sonata. A lyric. A sestina.
What are the “formal elements” in each? (goes well beyond Eaton, p. 80)
Cf. “Sonata Allegro form” – A(tonic) – B (dominant?) – repeat- Development, recapitulation A (subdominant?) – B tonic. Codas, bridges, etc. Development in bridges. Of bridges.
Fry, “the content doesn’t matter”(p. 81)
Clive Bell, religious feeling and “significant” form.
Liden and the poem.
c. connections between formalism and “special aesthetic emotions” (see Week IV, C 3.) Connections between formalism and aesthetic attitude theories. Connections to Plato, Kant.
NOTICE how fuzzy these distinctions actually are. One is always getting one KIND of content or another. There are lots of kinds. Eaton p. 83. Liden.
1. Kinds of context.
a. Conventional associations (iconlogicae!)
b. artist’s life.
c. art-historical facts
d. plain old historical facts
e. including, “social” facts (political etc.)
Cf. Steen. Homer. Durer.
Van Eyck (3.13). Liden. Bach’s
2. Marxist Aesthetics
a. Example: Berger on oil painting vs. e.g. ink. Possessions. The oils in English manors. (Cornwallis).
b. problem: exceptions Rembrandt’s biblical subjects. Other kinds of counterexamples. “Ownership” and “Three Beauties.” (Ukiyo-e – pictures of the fleeting world). Which comes first, (religious, metaphysical, ethical ) ideas or economic arrangements? Vaclav Havel.
3. Less sweeping contextualism. Walton.
a. Which “intrinsic features” get our attention depends upon “extrinsic” information.
b. in opposition to formalists – you cannot just stare (listen etc.) and notice certain “intrinsic” properties.
Examples (cf. 1 c): knowing what a sonnet is (historical facts) directs our attention to the 14 line, couplet–at-the end structure. Those are standard features which determine to what category a poem belongs.
The subject matter of the sonnet (love e.g.) is a variable feature. A poem with 13 lines could not be a sonnet, since having 13 lines is a “contra standard” feature so far as sonnets go (it disqualifies anything as a sonnet).
Other examples: how do you “hear” the opening line of Afternoon of a Faun? Which of its formal properties stand out? Must know the genre. Suppose you were expecting a classical symphony! (What, this is no classical symphony! What gives here?). Even knowing you have an orchestra makes a difference to the “properties” of that solo. (the radical of presentation).
Week IX Break
1. Institutional theory
a. Dickie; a work of art is
1. an artifact
2. some aspects of which have a “conferred status”
Cf. examples from Battin, avant garde art.
a. analogies – being married, being president. Observable features are not sufficient, or even necessary.
b. “the art world” does the conferring. Makes something a “candidate for appreciation.”
1. what are the boundaries of the art world?
2. HOW does
it confer (where is the “ritual”?).
3. what counts as “appreciation”? (Many different things). Give examples.
4. not just anything could be a candidate for appreciation. Examples;
5. don’t those who confer status do so for some reason? (Wolheim)
e.g. they think X is worthy of appreciation.
2. The Art World (Danto)
a. Art worlds are constituted by “institutions”
1. Danto goes beyond Dickie – there is no mystery about how certain objects get a status conferred
2. Works of art are constituted culturally or socially – like a language.
b. Identifying something as a work of art requires initiation into a culture.
Cf. Seeing scribbles and seeing a sentence (something “meaningful”).
3. Eaton – art and traditions of discourse.
a. x is a work of art iff x is discussed in such a way that audiences attend to intrinsic properties of x considered worthy of attention in aesthetic traditions (e.g. moulding in late Renaissance art).
b. worthy? What about “bad” art?
1. Both bad and good art invite attention to the same sorts of properties, i.e. those that are attended to within current aesthetic traditions.
2. In bad art, those properties do not please (don’t please whom?).
a. Certain basic forms or structures are found across cultures – cf. fairy or folk tales.
1. problem – James writes novels for sure – they fit that genre. But, THAT is not what matters most about them. What does?
a. The “elements” that constitute art works are not fixed by any interpretation or linked to anything permanent.
1. the elements (e.g. words, sentences) interact with each other, with larger contexts, endlessly. Cf. the sentence on p. 99, Eaton.
1. Texts (etc.) as occasions for “play” which never comes to an end.
2. Spy the “ideological” constructions and de construct them.
Representation? Resemblance? Convention? Iconic symbol?
Durer: Melencolia I
Denoting? Referring? Symbolizing? DEGREES of iconicity.
Reference within a system.
Portrait of Nicolaes Ruts
Oil on panel
118.5 x 88.5 cm
The Frick Collection,
Fan K'uan (d. after 1023)
Travelers Amid Streams and Mountains
Ink and color on silk
Franz Marc; Blue Horses
Winslow Homer “
Van Gogh - Shoes
Questions about interpretation:
Setting works from different periods “side by side.”
Mondrian Yellow, Red, White
Mondrian Broadway Boogie Woogie
Using Mondrian to illustrate a distinction between “style” and “repertoire” (Wollheim)
Setting an individual artist’s works “side by side.”
ART AND OTHER VALUES
Tolstoy’s (good) question. Either get a good answer or change your lives!!!
1. examples. The flower vase (p. 127)
a. economic value
b. sentimental value
c. historical value (cf.
d. art historical value
e. cultural value
f. ethical value
g. religious/ “philosophical” value (e.g. epistemic value)
1. c, e, f, g, as “truth” values (cf. art as “language”)
2. what does “aesthetic value” amount to when separated from 1a-g?
a. inherent value?(cf. formalism, theories of “aesthetic experience”, “disinterested contemplation” etc. Contrast with?)
b. consequential value? (cf. Plato –
Tolstoy – (think of what he DID value)
3. Art and ethical values.
N.B. Eaton connects consequential value in the arts with consequentialist ethics. THAT IS MOSTLY CONFUSING OR A MISTAKE.
Illustration: Plato is certainly not a consequentialist in ethics. But he evaluates art in terms of its value for knowledge. What makes the arts bad, or suspect, is that they make people even more ignorant than they already are (i.e. it has that bad consequence). But what is bad about being ignorant is not that IT has bad consequences, but because it is an intrinsically bad condition to be in.
One other sorry defect in Eaton: she begins by trying to limit the discussion of art/ethics to “consequentialist or deontological” ethics. But it seems pretty clear that the arts connect up with ethics largely through “VIRTUES” and some of her examples even illustrate that, e.g. Putnam, Richards.
Illustration: coming to see the world in various (richer?) ways might contribute to development of “character.” Virtues are “character” traits.
a. art and violation of moral principles (examples from Battin)
b. art and production of bad consequences (examples from Battin)
c. Pornographic “art” as an example.
1. moral principles
2. bad consequences
3. Consider that the bad consequence might be the “stunting of growth.”
4. Art and Public Policy
a. The NEA. Serrano.
b. Public art. Serra’s “Tilted arc” (v. Battin p. 181-193)
c. Public commemoration, history ( the
EATON on the AESTHETIC
1. aesthetic experience is –
2. aesthetic value is –
a. problem with 2. The notion of “pleasure” is either too broad or too narrow.
1. too broad if it includes any “pro attitude.”
2. too narrow if construed as mere sensuous enjoyment.
3. In fact, Eaton’s account still belongs to the “grand narrative.”
1. the (300 year old) GRAND NARRATIVE.
It consists of these ingredients:
a. art has now (by the 18th cent.) “freed” itself from service to anything extraneous to itself. And that is supposed to be progress! (Be able to illustrate the theoretical expressions of this idea from course content).
b. it is progress because when it is thus freed, art “comes into it own”, finds its proper place in life. Freed art can shape our engagement with it in the proper or appropriate way.
1. properly, it is the object of “engrossed contemplation.”
c. art releases us from social fragmentation, reintegrates human life into a “pristine” condition (the romanticist idea).
1. Art works as objects of transcendent worth.
2. The prophetic/messianic mission of art.
3. Art works as an avenue to “God”
2. Give up the grand narrative!!
a. art works, and artists, are not the unsullied items that the grand narrative supposes. (cf. Pound)
b. and anyway, art works are still immersed in “uses” of all sorts.
2. commemorative or public art
3. liturgical art
4. “political” art
c. the “mode of engagement” with such art is NOT disinterested, engrossed contemplation. (cf. kissing the Vietman memorial and kissing icons). Consider also the mode of engagement with certain photographs.
Lillegard (yah, that’s right)
1. A codicil to Wolterstorf. Art is not at an end. The Grand Narrative is?
Ordinary Narratives, tradition rooted in day to day contingencies, constitute art works in ways that make a great variety of modes of engagement possible.
a. Inta Ruka (narratives of birth, suffering, death, for remembering or memorial) Contrast to Adams, or Maplethorpe.
b. Gorecki – memorial, again. The virtues of remembering “roots.” History filtered consciously (rather than unconsciously!)
c. Rouault – memorial, protest, religion.
d. Kate Campbell (local history and…)
e. Anselm Kiefer (ironical /serious use of mythical, architectural, literary traditions, historical memory etc.)
Rouault – Old King - traditions of production, thought and belief, social protest, testimony to conscience, le sang du pauvre.
Two smaller works from the same period, principally in watercolor, combine darker themes with lush visual presentations. Your Golden Hair, Margarete (1980) quotes Romanian Jewish poet Paul Celan's poem Death Fugue, which is set in a concentration camp. Kiefer has sometimes depicted the German heroine's locks as straw adhered to the canvas; in this work, they appear in watercolor as sheaves of wheat in a field.
Winter Storm Clearing
In each of his images
Think of Bullough in the fog!
Mapplethorpe: formalism + eroticism…not exactly a new idea! Contrast Adams and Mapplethorpe to Ruka.
Think Georgia O’Keefe!
Georgia O’Keefe - Jack in the Pulpit II
"A phrase O'Keeffe used in a letter to Anita Pollitzer, "Tonight I walked into the sunset". (11 September 1916), is like all of her best art: immediate, concrete, all-encompassing, with a surprising syntax. Sunset is the time when the world appears least structured, when forms tend to dissolve and are replaced by new colors and sensations. O'Keeffe acted to suspend time, producing art that would capture the transient. For example, O'Keeffe made of a flower, with all its fragility, a permanent image without season, wilt, or decay. Enlarged and reconstructed in oil on canvas or pastel on paper, it is a vehicle for pure expression rather than an example of botanical illustration. In her art, fleeting effects of natural phenomena or personal emotion become symbols, permanent points of reference.” (Jack Cowart. Lillegard’s emphasis). Cf. Zen flower arranging, which stresses precisely the fact of contingency, impermanence, the approach of decay.
Where does art leave off, and decoration begin? You would never feel tempted to ask that question about Ruka’s work. It clearly is not decoration. The same cannot be said for O’Keefe (which is not to say that her art is not valuable in various ways).
1. According to Weitz it is not possible to define ‘art.’ Why?
2. Mention some reasons for thinking that what makes something a work of art is the intentions with which it was produced. Some reasons against. Give an example or two.
3. ditto 2, but “circumstances in which it was produced or presented.”
4. Eaton proposes a fourfold classification of aesthetic theories: they are
5. Eaton mentions six construals of ‘X (a work of art) expresses E.’ What are they? What is a point in favor of each one? A point against?
6. If artists express emotions in their works, does that make art “irrational?” Why would anyone think so? What would Plato have said? Aristotle?
7. Give some examples of the sorts of qualities or properties that a person of “taste” attends to in art works. Give a criticism of the view that certain qualities are noticed only by a person of taste, whereas others can be noticed by just about anyone. Use Sibley on Van Gogh for examples.
8. What is meant by the “aesthetic attitude?” Illustrate the idea from Twain and Bullough.
9. One problem with aesthetic attitude theory is that it applies equally well to ________ nature and ___ works. It completely leaves the artist’s own activity out of the account.
10. What is the problem of negative emotions? Give some examples of art works involving negative emotions, and give an account of them that explains why people would be interested in such works.
11. Discuss two versions of the view that art works are the result of a special kind of creative activity.
12. What might be some examples of special aesthetic emotions?
13. Describe a case of an ordinary emotion that is included in a response to an art work, but in a special way.
14. Give an illustration of Dickie’s view that what distinguishes aesthetic response is the things that are attended to, rather than the manner of the response.
15. The following are ways in which art works, or elements in them, have been thought to refer to something: they might refer by representing; imitating; resembling; being conventionally related to; standing for. Illustrate each of these using the following examples (a single example might illustrate more than one of these); Handel’s ‘All we like Sheep;’ The bat, and the compass, in Durer’s Melecholia I ; Hobo signs and road signs; a musical phrase (Handel); Brancusi’s “The Lover’s kiss;” The figures in Jan Steen’s painting; Battin 3.13, 3.15.
16. Discuss two ways a work of fiction might be “true.” Give examples.
17. Discuss two ways a painting or other visual art might be “true” Give examples.
18. Discuss ways in which fiction or visual art might be “false” or “lie.”
19. Various “signs” or symbols in art works vary in their degree of “iconicity;” Illustrate with reference to a piece of music, a painting or sculpture. You can use the examples from question 15 if necessary.
20. Could works of visual art, or of music, have a “syntax?” Illustrate.
21. Give a few examples from painting, music, literature, of “formal” properties.
22. Formalism: use several examples of different sorts of art works that might suggest different answers to this question: is it possible to notice only the formal properties of an art work?
23. Give some examples of art works that exist only as types. As tokens.
24. Formalists supposedly claim that the only properties of works of art that deserve attention are those immediately present in the work. Critique this notion by connecting it to
a. the “grand narrative”
b. Clive Bell’s “pantheist” religiosity
25. Mention as many different kinds of “context” of an art work that might be important to its status and significance as an art work as you can think of.
26. Give an illustration of Marxist contextualism using two of the works above.
27. Give several illustrations of Walton’s anti-formalist notion that which intrinsic features of an art work get our attention depends upon extrinsic information.
28. Explain what Dickie means when he claims that art works have features due to a “conferred status.” Then offer three criticisms of the institutional theory.
29. Give reasons for thinking that there is no such thing as pure “factual” description of an art work.
30. Give reasons for thinking interpretation of art works is not merely “subjective.” Include a discussion of the kinds of “reason giving” involved in interpretation, and the idea of the non-transferability of critical judgments.
31. Use the cited interpretations of Winslow Homer and Gaugin (spirit of the dead watching) to illustrate “political correctness” in criticism.
32. Illustrate the idea that criticism consists in “inviting and pointing” by considering the following cases:
a. pointing out art-historical facts by setting Durer and Ostade side by side (what fact become noticeable?) Give another example.
b. pointing out facts about repertoire by putting one artist’s work next to another of that same artist’s work (cf.Mondrian).
c. pointing out various “unities” in a Tolstoy story.
d. pointing out the fact that the early 18th
cent. is an age of increasing “individualism” (cf.
e. pointing out
33. For each of 32 a-d, connect to question 25.
34. Illustrate Eaton’s definition of art by reference to typical reactions to Egyptian wall paintings on the one hand, a Rembrandt portrait on the other.
35. Show how interpretation of an art work can be improved or made more correct by attending to
a. the artist’s intention
b. art-historical facts
c. ordinary historical facts
d. facts intrinsic to the work.
Illustrate each of these from Battin, Eaton or class examples.
36. a.Mention at least three “other values” that might conflict with so-called aesthetic value. Illustrate (from Battin, Eaton, class). Begin with a discussion of Tolstoy, and mention Plato, and Riefenstahl.
b. mention at least three “other values” that might be enhanced by or communicated through art works. Illustrate. Mention Aristotle in one illustration.
37. Discuss several difficulties in separating “esthetic” value from other values. Illustrate.
38. Discuss issues of art and public policy in relation to Serrano, Maplethorpe, Serra. Sketch several positions (e.g. “there should be no public art” or “art should be constrained by moral norms”) and argue pro and con.
39. Describe the main ingredients in the “grand narrative” as Wolterstorff conceives it. Then illustrate some alternative (non-grand narrative) approaches to art works.
40. Illustrate each of the following (discussed in Wolterstorff and Lillegard) using examples from class, and point out different kinds of engagement with art works which these examples suggest;
a. art and memorial (give several)
b. art and narrative (esp. keep in mind ordinary life narratives, “birth, copulation, death” ).
c. art and morality (discuss virtues in particular. See quote 1 below)
d. art and religion
e. art and political statement
f. the end of art, decoration, and formalism.
g. art and testimony (testifying). To what? Truisms.
Include a discussion of the photography of Adams, Maplethorpe, Ruka wherever they fit.
photographs were the core of
1. A quote from “Spirit and the End of Art”
“A story, if it is going to work as a story and not just as an incidental means to the expression of an ideology, must show the unfolding of life through time. And that in turn requires continuity in its subjects. Now virtues are precisely such formed traits as provide continuity in a life. Our ability to appreciate a life as virtuous depends upon our grasp of the way in which existence gets form through passionate commitment and formed response. Virtues are precisely patterns of formed response, dispositions of thought, feeling and action which continue and enforce a coherent life. Vices too take on significance in relationship to possibilities of achieved continuity. We understand what it means for a life to fall apart because we understand what it means for a life to come together”
Another quote to reflect on;
2. “The photographs of Inta Ruka should be viewed as an amalgamation of images and text. . .it is a broad and clear manifestation of an ethical position. To quote Prof. Klavins; “it reconciles with reality.” Even though people may be inscribed in some particular social context. . .for me, you and her country people there is only one life with a beginning and an end.” (Helena Demakova, in Inta Ruka: My Country People)
3. “ In Rilke's poem "Death Experienced," in memory of
Countess Schwerin, he writes: "when, though, you went, there broke upon
this scene/a shining segment of realities/ in at the crack you
disappeared through: green/of real green, real sunshine, real trees."3
It is remarkable how close in spirit, even in detail, Rilke's lines are
to the following lines from a poem by Orville Kelly written to his wife
as he faced certain death from cancer: "Summer, and I never knew a
bird/could sing so sweet and clear,.. .I never knew the sky could be so
deep a blue, /Until I knew I could not grow old with you..."4 These are
characteristic expressions of the grace-filled enhancement of existence
which can only be realized in the face of my own death, fully accepted.
At this point I will enter a few remarks on cloning, as promised above. My original argument was intended to establish the claim that what we can sensibly say and think about ourselves depends in part upon the fact that we are the kind of beings who are begotten, born, reproduce in a certain way and die. Being born contrasts with being manufactured. We count as "our kind" others brought into existence in like manner, that is, through acts of begetting which are fundamentally biological and not subject to perfect control. Thus it is natural even for many people who are not religious believers to think of a child as in some way a "gift." These facts have, I believe, a direct relevance to what we should think about cloning of humans. In his testimony before congress the theologian/ethicist Gilbert Meilander argued that the idea of "begetting" is essentially opposed to "making" and that the fact that our children are begotten rather than made is fundamental to the conceptual background required for thinking of them as free and sharers in a common dignity, as more than just functions of our individual projects or outputs of our desires and wishes.(Meilander's statement, presented to the National Bioethics Advisory Commission on March l3th, 1997, is printed in First Things, #74 (July/August 1997), p.41-43.) The fear that cloning would alter that background is I believe well founded.
I have accused Sennett of gnosticizing tendencies. The gnostics' distaste for earthly life is rooted in a hatred of contingency, the limitations of space, time, mortal embodiment and our inability to control our lives, which seem to them obstacles to
"spiritual" release and realization. . .”
(From Lillegard, “Reply to Sennett and Wildman” in CrossCurrents 1998) |
With the Collapse of Royalist control in the British Isles and the uprooting of the rich and a great deal of the middle classes, the integration of British political elites in the Canadian political sphere has been characterized by its own unique problems.
THE CANADIAN EXPERIENCE IMMEDIATELY AFTER THE WAR:
When Canadian troops returned to their country after the conclusion of the peace with honour treaty, Canadian Prime Minister Robert Laird Borden received a marching column of troops headed by the tenacious and well-beloved Canadian general Sir Arthur Currie GCMF, KCB. Currie, esteemed heavily by both the Allies and the Central Powers for his brilliant leadership of the Canadian Corps, had quickly established a reputation for the Canadians as some of the most effective troops on the Western Front.
PM Borden, dogged and tired after holding a Union Government together for the duration of the war, appeared visibly strained and prematurely aged as he proclaimed: "I welcome you home from the fields of war, from which you knew no defeat!" Canadian attitudes to the war were divided. While Patriots and Imperialists took comfort some comfort in the fact that the Empire had not been defeated per se, the fact that the war ended in a unfavourable "draw" heavily in the German favour kindled great resentment, espicially in Quebec where conscription riots and opposition to the war seemed to be justified in the light of German dominance of the European continent.
In the immediate aftermath of the war, all the Union liberals who had not left the government quickly rejoined their fellow party members in the opposition and the government fell apart, divided as to how to take Canada into the future. Former liberal PM Sir Wilfrid Laurier's optimistic prediction that the 20th century would be Canada's century seemed laughable now to most Canadians. Laurier himself had passed away during the war, seemingly taking with him to the grave the sunny prosperity that had affected Canada during his leadership in the late 19th century.
Shortly after calling the election, PM Borden lay in his room at the Chateau Laurier (named in honour of the former liberal PM) and passed away. His last words were of regret that he had never returned to his native Nova Scotia to rejoin the Celtic Diaspora.
With many of the heavyweights of the old 19th century and pre-war political traditions dead, Canada weaved in and out of political obscurity and a deep sense of national malaise. All of this would be grimly impacted by the loss of the Home Islands.
With the loss of the Home Isles, the British liberals, still tarred with having lost the war under the leadership of Herbert Asquith and David Lloyd George (at least in Canada), the liberal tradition of Britain did not impact the Canadian liberal party as much as some thought it would. William Lyon MacKenzie King, who became liberal leader in 1919, had previously been Minister of Labour under Laurier and came from an interesting family. His maternal grandfather had been the leader of the Upper Canadian rebellion and had died in the United States in exile. After serving under Wilfird Laurier as Labour minister for a brief time in 1907, King took time to travel the United States and wrote extensively on Labour relations, becoming well acquainted with many leading American industrialists such as the Rockerfellers while doing his best to aid them in successfully resolving their labour disputes. He returned to Canada at the outbreak of the Weltkrieg.
With the Conservatives led by Arther Meighan after the death of Borden, King has met with mixed success in rallying the forces of Canadian liberalism. His greatest success has been painting Meighan as a yes-man to the British lords and exiles blamed for losing the Empire, but his policy of urging a peaceful foreign policy based on mutually dependent, integrated economies with other the nations within “The Anglo-Saxon World” (as he calls it) has antagonized the British and Canadian military elites and has proven difficult to reconcile with Canada’s massive fleet and imperial obligations in Delhi and the Caribbean. One of King’s highest priority both as Opposition Leader and Prime Minister seems to be helping the traditional political parties in the USA retain control as part of the Anglo-Saxon World. It is rumoured that if civil war comes to the USA, King may side with the America First movement if they legitimately take power, if only to prevent the growth of syndicalism in North America…
An economic upturn and expansion of industrialization was brought about by the infusion of British gold and technology brought by the British exiles but it has not brought a corresponding rise in worker’s rights due to suspicions of Syndicalism. Left wing elements in the Prairies organized a general strike in Winnipeg which was brutally repressed by the Royal Canadian Mounted Police, egged on by the British and Canadian elites who desired to stamp out any possibility of another revolution occurring. However the extent that many of the RCMP officers went to shocked many within the Canadian body politic (it was rumoured that the RCMP has been infiltrated by a group of men who are the sons of former members of the old Special Branch), and the Liberals under King reached out to members of the Canadian Commonwealth Federation (CCF) declaring that they were merely ‘Liberals in a hurry.’ King’s charm and reaching out to the left-wing in order to form a united opposition to the Conservatives has resulted in Canadian socialism and liberalism adopting a policy of gradual, legal reform. With this alliance, King was able to win the next federal election and regain the Premiership.
This has not endeared King to many conservatives in Canada. While King’s policies are clearly constructed in view of how much the Canadian people will tolerate, his oppurtunism and accumulation of political alliances among the respectable members of the Left has built a powerful counterpoint to conservatives such as Arthur Meighen and R.B. Bennett. As dangerous as it seems, King is no fool. While his eccentricities are well known, (his doting love for his Scottish Terrier Pat, his interest in séances and the occult and his devotion to his mother) he maintains an iron grip on his party’s discipline and arranges for himself to carry a safe riding every election.
Relations with the Liberals and the Royal family are strained. The King has no time for the Liberal leader, leading wags in the press to talk about the war between the King’s of Canada. The Prince of Wales, the darling of Canada’s media, finds King to be nothing but a mere colonial and spends most of his time putting in valuable facetime for the Monarchy in Canada, especially in Quebec where he is doing his best to woe the French Canadians with his charm (rumours are swirling that he is involved with a Roman Catholic Frenchwoman, a commoner no less, and may in fact cost him the crown). His brother Albert however, possessed of a powerful sense of duty, has granted King a great deal of mutual respect, if not public endorsement, especially in regards to King’s contention that each British Dominion is a separate state with equal rights within a larger Empire, each with its separate crown for its head of state.
The other sons of King George V, HRH Henry Duke of Gloucester and HRH George Duke of Kent have kept a low profile in the Canadian public life, though the Duke of Kent’s interest in the development of airpower has lead to his active engagement and elevation within the Royal Canadian Airforce. He has been instrumental in establishing the Empire Air Training Program in the Province of Alberta where pilots from all over the Empire, from the Pacific, South Africa (though less and less are coming from this troubled part of the Empire), the Caribbean and even some from India.
Henry has remained in the army and patronizes the Royal Military College in Kingston as an honorary colonel. It is suspected that he will be dispatched to one of the other Dominions as Governor-General to increase the connection between the Crown and its subjects.
CANADIAN ELECTIONS FROM 1917
1917 - Conservatives, led by Borden, are re-elected with a majority as part of a pro-conscription unionists coalition, which had former Liberals and Conservatives in the cabinet. The Unionists defeat Laurier's anti-conscription Liberals in the most bitter campaign in Canadian history.
1922 - Liberals, led by William Lyon Mackenzie King, win a minority government, defeating Conservative prime minister Arthur Meighen. The Conservatives are reduced to third place in the House as many Center Left wing MP’s win elections in the West and British Columbia under the new Progressive party. However, the Progressives decline the title of official opposition, leaving open the possibility of a coalition with the Liberals. Meighen becomes opposition leader.
1925 - Mackenzie King's Liberals hold on to power with the help of Progressive Robert Forke, despite Conservatives, led by Arthur Meighen, winning more seats. The Revolution in Britain occurs and the Union of Britain is proclaimed. A political crisis ensues as a flood of British political refugees pours into Canada. The Progressives, painted with the brush of Syndicalism are viciously attacked by the Conservatives in the media. Some street fighting is reported between members of the Right and Left wing militias. The Royal Canadian Mounted Police restores order and Prime Minister King declares a National Emergency and invokes the War Measures Act. The Royal Navy re-establishes itself in Canada as the Union of Britain consolidates its control of the Home Isles. Without a centralized Imperial government, the Empire is thrown into confusion, the Royal Navy, acting mostly of its own initiative is able to maintain control of the Caribbean. The British Indian Army re-groups in Delhi and salvages as much of the former Raj as it can. Australia and New Zealand, without the protection of the Royal Navy (later to become a point of great contention between Ottawa and Canberra) merge into a political union. The German Empire seizes most of the Strategic naval stations of the British Empire, allegedly to prevent a world crisis, but it becomes clear as most of the Empire plunges into rebellion and civil war that the British fleet cannot insure order and the German presence becomes permanent. The declaration of a National emergency plays into King’s hands as he is now able to stifle accusations of corruption and patronage within his government (that actually turned out to have much validity) and the more centrist of the Progressives jump ship to the Liberal party which wraps itself in the flag to weather the Crisis. With a stronger Liberal party and a emasculated Progressive party, the Conservatives bide their time, constantly criticizing King’s failure to secure more of the Empire (though realistically, the tiny Canadian army was only able to send a token force to the Caribbean). King, realizing that the Empire is lost but that the admission of such a fact is political suicide, embarks upon a new Imperial Policy as the most powerful British state left standing. The presence of a British government in exile is a threat to Canadian unity and King’s power, and under the provisions of the War Measures Act and Martial Law, King promulgates a legislative union act, merging the two governments into one under his control. He is able to outmanuver and replace all British leaders who could conceivably try to become PM of Canada in the highly charged emotional mindset of the fall of Britain. After asserting his control over the Cabinet and parliament and making sure that only British MP’s who will follow his line and get elected in their own right in Canadian ridings, King moves on to the armed forces. Seeing the Royal Navy as the only common institution of the Empire still intact, the Fleet is ordered to divide up and deploy squadrons to Karachi, Australasia and South Africa. This does much to reassure the Australasians and South Africans, though it won’t be permanent.
1930 - Conservatives, led by R.B. Bennett, win a majority, defeating the Liberals under Mackenzie King. With the influx of many British Conservatives such as Stanley Baldwin, Neville Chamberlain and Winston Churchill, the Conservatives have been able to present themselves as the party of Imperial unity. The endorsement of King George V, blatantly meddling in Canada’s political affairs seals the election for Bennett whose cabinet is 100% Canadian. However, Bennett relies on the financial, political and philosophical guidance of British Exiles who exist in a ‘shadow cabinet’ of Canada’s traditional clubs for the Elite such as the Empire Club in Toronto, the Orange Order of Ontario and the Chateau Clique in Montreal. Even under the Liberal government of King, the Army and Navy staffs were packed with many of Britain’s top generals and admirals. Though a few generals such as Vanier (the top ranking French Canadian in the army), Crerar and McNaughton have attained positions of influence in the army, the Navy is firmly in the hands of British admirals. The Air Force is the only Canadian institution to remain almost equally divided between Canadians and British Exiles. Proclaiming that the ultimate goal of Canada is the reclaimation of Britain, Bennett recalls much of the British Fleet from the Pacific and the entire South African station, causing great resentment in the other Dominions and wages a inconclusive (though popular) war of skirmishes with the Syndicalist Navy in the Atlantic. While the Navy seems well prepared to renew the war against the Union of Britain, the Canadian army remains small and is in need or reorganization if it is to become capable of forcing a landing in Britain. Many Canadian officers look enviously at the Australasian army, which is no more advanced, but is larger, better organized and better led as several of the more progressive British generals settled in Australasia.
1935 - Liberals, led by Mackenzie King, defeat Bennett's Conservatives with a majority. The election was seen mainly as a referendum on Bennett’s Imperial Policy. While it played well at home to wage a tit for tat naval war in the Atlantic for a about a year or so, the public has become weary of a struggle with no conclusive end or objective in sight. The Australasians, whose own economic expansion brought about by its own influx of British Exiles and currency (though much less than the one that affected Canada) has stalled and continues to resent the loss of the majority of the Australasian station of the Royal Navy. Many in Canada fear that the larger Australasian army may not march for the Empire if it means leaving their shores undefended. King’s sensitive balancing of the old Entente Alliance to keep the Dominions firmly within an Global Anglo-Saxon Bloc has resulted in many tense negotiations between Ottawa and Canberra to re-deploy some of the Royal Navy back to the South Pacific, especially in light of the expanding Japanese Navy, which is still technically an “ally”.
FOREIGN POLICY AND DIPLOMATIC RELATIONS WITH THE OLD EMPIRE STATES
Abroad, the old Boer nemesis in South Africa has resurfaced. South Africa was the one Dominion of the Old Empire to receive only a trickle of Exiles and the share of the spoils they were able to save from Britain was miniscule. Even King may not be able to bring the South African’s back into the Empire, though he hopes that they will at least remain as an associated power and maintain the preferential trading agreements that see many of South Africa’s raw materials head to feed Canada’s large industrial appetite.
The British Naval Station in Delhi remains strong enough to be threatening, but too weak to maintain Order in the Indian Ocean and is growing old in the face of the German Naval Squadron based at Ceylon, the Canadians maintain it to maintain a link with the aging British technocrats in the Delhi government.
The fall of the Old Empire and the rebellion of several subject peoples has resulted in a great deal of racism in Canada and Australasia against people who are visibly Non-British, and the Conservative party, back when it was still in power, urged the Australasians to deploy a sizable portion of its army to Delhi to “keep the natives in line”. This further antagonized the rocky relationship between Canberra and Ottawa, and King astutely dropped the matter. While they are not Anglo-Saxon, King has made a point of declaring on no less than 7 public occasions that India has “been a part of the Empire longer than most of Canada and that there exists a special place for them, side by side with their fellow citizens of the Empire…” The fact that King used the word citizens, as opposed to subjects has caused a great deal of debate over the future of the Empire in political circles. With the advent of more sophisticated communication technology, many proponents of a global federation of Imperial British states is gaining more prominence within Liberal circles, especially as it would place a reclaimed Britain on an equal footing with Canada within such a political arrangement.
In Quebec Maurice Duplessis, as leader of the paternal-authoritarian-like Union Nationale, stands unopposed in provincial politics and successfully eliminated overt support syndicalism and ensuring the survival of the Catholic identity of Quebec. His party has used the election slogan "Survival" in every provincial election to paint Quebecois self-image as a religious struggle for God.
Duplessis however has been more than willing to let Anglo companies build branch plants within Quebec and build up an impressive commercial presence so long as they did not allow unions and did not try to intervene in politics at all. He has even allowed some of them to have English-only workplaces. He maintains cheerful and amicable relations with the corrupt and patronage driven government of Mitch Hepburn, the scandalous "Liberal" premier of Ontario.
In light of the fall of the Empire, Duplessis has enhanced relations with the Italian federation and its dominant Catholic political culture, especially in hopes that a weak pope will be elected when the current incumbent dies. If such an election took place and the Pontiff would not interfere with Quebec, Duplessis is eager to maintain his close alliance with the Catholic Church. The Church essentially runs education and social welfare as a result of this arrangement and very few Quebecois receive a higher education. Duplessis also maintains a relationship with National France, not because of a kinship to the European French (whom many Quebecois accuse of abandoning them to the English), but to stamp out any possible Syndicalist infiltration by Commune agents.
Assuming a strong Pope like Innitzer gets elected... Duplessis may be forced to reconsider his alliance as he would not want the Holy Father looking over his shoulder. King has long given up on trying to engineer a liberal government for the Province of Quebec and both Conservative and Liberal governments have refrained from interfering in Quebec politics so long as Duplessis does not stoke disloyalty or interferes with Federal politics.
While R.B Bennet and the British elites wanted to also recapture Ireland, King has embarked on a policy of reconciliation with the “Lost Dominion.” In the event of war, King have suggested sending a diplomatic offer to Michael Collins that offers a public, binding, international decree by the British crown to recognize the sovereignty of the Irish Republic and a vow to never violate that sovereignty, and in return for port rights and airstrips in Ireland to carry out the war against the UoB.
(Amazing Canadian Backstory by UnitedEmpireLoyalist, Cheers)
1936, The Present
Ever since the Proclamation of the Union of Britain, King George V has been plagued with sickness, a sickness which has intensified in recent months. The fate of both the King and of the Empire is unsure, both have begun to slowly wither and die. Next in line for the throne is Edward, Prince of Wales, while R.B. Bennet and the Exiles are waiting for their chance to use the Prince to their advantage. No matter what happens in the first few months of 1936, God Save the King!
- Kaiserreich: Legacy of the Weltkrieg 1.0
- Settings on Normal/Normal
Kaiserreich: Legacy of the Weltkrieg is a mod for HoI2 Doomsday and Armageddon which asks the question: "What if the Germans won WWII?" As you have read, Revolutions at home have forced the British Royal Family to flee to Canada and for the French Nationalists to cross the Mediterranean Sea to Africa.
Now, a decade later, Berlin's economy is in recession and political divisions run deep in the United States. General Strikes by the IWW bring the economies of Chicago, Detroit and Green Bay to a halt while the "American First Council" agruably has more power in the Midwest and Texas than the Federal Government while the West Coast of the United States prospers.
The Commune of France is the leader of a Worldwide radical movement whose goal is to bring down the Reactionary governments in Europe and abroad and establish rule by Trade Unions across the World.
In Russia, political dissent has reached its boiling point and divisions are begin to appear in Spain while the Ausgleich, the Re-Negotiation of the Austro-Hungarian Comprimise is nearing closer.
The economies of South America are dependent on Berlin, its anyones bet what would happen to the Continent if the recession in Germany was to stop sliding and to plunge.
Success or Failure; Whether London is retake or the Revolution follows the Royal Family across the Atlantic, I hope you'll be entertained. |
As Internet Protocol (IP)-based networks are increasingly deployed, packet-based applications such as voice over IP (VoIP), IP video, video telephony (VT), integrated voice and email, and instant messaging (IM), have emerged. The ability to integrate these services over the same network has become more important as customers appreciate and demand the bundling of multiple services. Many of these services benefit from session based connections between communicating network devices. For example, rather than having each data transmission between two devices be considered independent from each other, a series of data transmissions may be logically grouped together as a session. As session based traffic increases in the network, the problems of how to provide redundancy and load balancing among a cluster of session handling servers have to be addressed.
BRIEF DESCRIPTION OF THE DRAWINGS
FIG. 1 a is a diagram of an exemplary system for implementing a congestion control method;
FIG. 1b is a system diagram of an exemplary upstream network device;
FIG. 2a is a representation of server and cluster capacities;
FIG. 2b is a representation of server and cluster capacities in the event of a server failure;
FIG. 3 is a flowchart depicting exemplary steps and decisions related to a process for managing congestion control and server redundancy;
FIG. 4 is a flowchart depicting exemplary steps and decisions related to a process for updating server utilization factors; and
FIG. 5 is a flowchart depicting exemplary steps and decisions related to a process for scheduling packet data traffic to servers in the cluster.
Exemplary illustrations of a congestion control method for session based network traffic are described below. In the interest of clarity, not all features of an actual implementation are described in this specification. It will of course be appreciated that in the development of any such actual illustration, numerous implementation-specific decisions must be made to achieve the specific goals of the developer, such as compliance with system-related and business-related constraints that will vary from one implementation to another. Moreover, it will be appreciated that such a development effort might be complex and time-consuming, but would nevertheless be a routine undertaking for those having the benefit of this disclosure.
Referring now to the drawings wherein like numerals indicate like or corresponding parts throughout the several views, exemplary illustrations are provided.
FIG. 1a illustrates a system 100 for use with a congestion control method for session based network traffic. Details of the elements depicted in the figures are included following a brief functional overview of the system 100 and method. Packetized data on a packet data network 105 may be transmitted to and from any of a number of interconnected network devices. For example, end user devices such as a Personal Computer (PC) 110, a mobile phone 115, a television 120 via a set-top-box 125, etc., may be the recipient or initiator of data traffic over the packet network 105. Servers 130a-c may receive and process data traffic from the end user devices. Moreover, the servers 130a-c may provide responsive data traffic to the end user devices. When multiple servers 130a-c collectively handle the same type of requests and data traffic, they may be arranged into a cluster 135.
When arranged as the cluster 135, additional measures may need to be implemented to ensure that the servers 130 are utilized in a suitable manner. For example, it may be desirable to have similarly situated servers 130 handle corresponding amounts of data. Accordingly, a load-balancing scheme may be implemented to assign certain amounts of data processing to particular servers 130 of the cluster. Additionally, certain amounts of the data processing capacities of the servers 130 may need to be reserved for redundancy. A redundancy scheme may allow the cluster 135 to maintain normal operations even with a failure of one of the servers 130.
By focusing on the capacity of each server to process data as well as on the amount of data actually processed by each server 130, a single cluster manager may be able to implement both redundancy planning and load-balancing schemes. A router 140, or any other network device acting as the cluster manager, may have the responsibility of tracking the amount of data actually processed by each server 130 and scheduling new data traffic to a particular server. Accordingly, the router 140 may provide a congestion control module 145 to implement the redundancy planning and load-balancing schemes.
The packet network 105 may be a packet switched communication network such as an Internet Protocol (IP) network. The packet network 105 generally interconnects various computing devices and the like through a common communication protocol, e.g. the Internet Protocol. Interconnections in and with the packet network 105 may be made by various media including wires, radio frequency transmissions, and optical cables. Other devices connecting to and included with the packet network 105, e.g., switches, routers, etc., are omitted for simplicity of illustration in FIG. 1. The packet network 105 may interface with an IP Multimedia Subsystem (IMS), which integrates voice, video, and data communications on the same network infrastructure.
The PC 110, mobile phone 115, and television 120 by way of a set-top-box 125 merely represent three of many possible network devices capable of connecting to the cluster 135 for session based data processing. Each of these devices has computer software including an implementation of the network protocol needed to communicate over the packet network 105. Additionally, the devices may also implement higher level protocols to interface with the session handling facilities of the servers 130. For example, the mobile phone 115 may include instructions for conducting a voice communication based session over the packet network 105. Likewise, the set-top-box 125 may include instructions for conducting a video based session over the packet network 105.
The servers 130 of the cluster 135 generally provide session handling capabilities. For example, they may be able to initiate a session based on a request from one of the network devices 110, 115, 125, etc., process data during the session, and terminate the session as necessary. Session Initiation Protocol (SIP) is a signaling protocol for initiating, managing and terminating application sessions in IP networks, and provides a mechanism to allow voice, video, and data to be integrated over the same network. Accordingly, the servers 130 may implement SIP for session handling.
The router 140 may interconnect one or more computing devices, e.g., servers 130, to the packet network 105. Moreover, the router 108 may establish and operate a local area network (LAN) for the servers 130, and may route certain communications thereof. For example the computing devices 110 may be connected to the router 108 using a wireless connection, a network cable such as a “Cat5” cable, or the like.
The router 140 may further act as the cluster manager of the cluster 135. The cluster manager may be an upstream network device to the cluster 135 such as a device positioned between the cluster and the packet network 105. However, in another exemplary approach, one of the servers, e.g., 130a, may be designated as the cluster manager. In such an approach, the designated server 130a may be considered upstream in that it acts as a gateway of the cluster 135 for incoming network traffic.
The congestion control module 145 may include computer instructions for implementing redundancy and load balancing schemes. The congestion control module will be provided by the device acting as the cluster manager, e.g., the router 140 as depicted in the exemplary approach of FIG. 1. As will be discussed in detail below, the congestion control module 145 may continuously schedule incoming traffic to particular servers 130, calculate utilization rates of the servers, monitor the cluster 135 for server failures, and update redundancy factors based on the state of the cluster. Positioning the congestion control module 145 upstream from the servers 130 may allow for the estimation of current utilization rates of the servers. Estimating utilization rates by the congestion control module 145 may eliminate time consuming bi-direction communication between the servers 130 and router 140 to determine the actual current utilization rates.
FIG. 1b illustrates the elements of an exemplary upstream network device, such as the router 140. As illustrated, the router 140 may include elements for guaranteeing the quality of service for different classifications of data. For example, the incoming packet data 155 may first encounter a classifier 160. The classifier 160 may inspect the header of the packet data 155 to identify the proper classification. A marker 165 may write, or rewrite, a portion of the packet header to more clearly identify the determined classification. For example, the marker 165 may write to the Differentiated Service Code Point (DSCP) field of the header.
A meter 170 may track the amount of incoming data 155 processed by the router 140. Moreover, the meter 170 may track the amount of data 155 for each classification. A policer/shaper 175 may use the tracked amounts from the meter 170 to enforce particular traffic routing polices, e.g., quality of service guarantees, service level agreements, etc. To enforce a policy, the policer/shaper 175 may drop packets if the tracked amount of data 155 exceeds the service level agreement. Additionally, the policer/shaper 175 may buffer or delay traffic that fails to conform to policy being enforced. A scheduler 180 has the responsibility of deciding which packets to forward to the cluster 135 for data processing. The scheduler typically bases its decision on the priority levels of the data as well as on the service level agreements. However, the scheduler 180 may be further influenced by the congestion control module 145. In one exemplary approach, the congestion control module 145 may be integrated into the scheduler 180. However, in another exemplary approach, the congestion control module 145 may be distinct from the scheduler 180 while still providing input thereto.
FIG. 2a illustrates a representation 200 of server and cluster capacities during normal operation. Each server 130 may have an actual capacity 205, which represents the normal or default ability of a server to process data. As depicted, each server 130 has the same actual capacity 205a-d. However, in another exemplary approach, servers 130 with different actual capacities may be part of the same cluster 135.
The actual capacity 205 may be artificially limited or reduced to an available capacity 210. The reduction of the actual capacity 205 provides redundant capacity 215 that is reserved for periods in which there is a server failure. The reduction of the actual capacity 205 of each server 130 may differ. Accordingly, a respective redundancy factor, ri, i=1, 2, . . . , n, where 0≦ri≦1, may be established for each server 130. The redundancy factor ri states the amount of redundant capacity as a fraction or percentage of the actual capacity 205. Accordingly, the available capacity 210a for each server 130 will be expressed as (Actual Capacity*(1−ri)), while the redundant capacity 215 will be expressed as (Actual Capacity*ri)
If ri=0, it is assumed that the server 130 is expected to have an available capacity 210 equal to the actual capacity 205. Accordingly, any server with a redundancy factor of zero (ri=0) will not provide any redundancy to the cluster 135. On the other hand, if ri=1, then the server will limit its entire actual capacity 205 for redundant capacity 215. Moreover, any server with a redundancy factor of one (ri=1) will not handle any data processing requests unless there has been a failure of another server in the cluster 135. In general, higher redundancy factor values limit larger amounts of actual capacity 205 for redundancy capacity 215.
The available capacity 210 may be used to handle data processing sessions from remote devices. A current usage 220a-d of each server reflects the amount of the available capacity 210 that is currently being used to handle data processing sessions. The current usage 220 will typically fluctuate as the server 130 receives new sessions and completes others. Moreover, the data processing demands on the server 130 may vary throughout a session. The router 140, or equivalent cluster manager, may determine which server 130 should handle newly received packet data. As will be discussed in detail with respect to FIG. 4, a utilization factor may be defined as the ratio of the estimate of current usage 200 over the available capacity 210. In one exemplary approach, the router 140 may implement packet data scheduling decisions by directing incoming traffic to the server meeting at least one threshold criterion. In one illustrative approach the criterion includes a consideration of utilization factor and in another approach the criterion is the lowest utilization factor.
Because the servers 130 act together as a cluster 135, the sum of the actual capacity 205 of each server defines an actual cluster capacity 240. Similarly, the sum of the available capacities 210 of each server defines an available cluster capacity 250. The sum of the redundant capacities 215 of each server defines a cluster redundant capacity 255. The actual cluster capacity 240 will remain constant so long as there are no server failures in the cluster 135. Likewise, the available cluster capacity 250 and the cluster redundant capacity 255 will remain constant so long as there are no changes to the any of the redundancy factors (ri). However, a cluster usage 260 will fluctuate as the sum of the current usage 220 of each server varies.
FIG. 2b illustrates a representation 202 of server and cluster capacities during a failure of one of the servers. As indicated by the X, the actual capacity 205a is currently unavailable due to a server failure. Accordingly, the actual cluster capacity 240 is reduced by the actual capacity 205a of the failed server 130. Because of the failure, the redundancy capacity 215 of the remaining servers may be reallocated as available capacity 210. Moreover, the redundant factor of the remaining servers may be set to zero (ri=0) in order to cause the available capacity 210 to fully encompass the actual capacity 205. The current usage 220a of the failed server 130 represents sessions with incomplete or unfinished data processing. Accordingly, the sessions encompassing the current usage 220a may be redistributed to the remaining servers of the cluster 135, thereby increasing the current usage 220b-d levels thereof. As expected, the cluster usage 260 will typically remain unchanged.
The router 140 and servers 130 may be any general purpose computing device, such as a PC, or a specialized network device. The router 140 and servers 130 may have software, such as an operating system with low-level driver software, and the like, for receiving signals over network links. The operating system may also include a network protocol stack, for establishing and accepting network connections from remote devices.
The router 140 and servers 130 may employ any of a number of user-level and embedded operating systems known to those skilled in the art, including, but by no means limited to, known versions and/or varieties of the Microsoft Windows® operating system, the Unix operating system (e.g., the Solaris® operating system distributed by Sun Microsystems of Menlo Park, Calif.), the AIX UNIX operating system distributed by International Business Machines of Armonk, N.Y., and the Linux operating system. Computing devices may include any one of a number of computing devices known to those skilled in the art, including, without limitation, a computer workstation, a desktop, notebook, laptop, or handheld computer, or some other computing device known to those skilled in the art.
The router 140 and servers 130 may include instructions executable by one or more processing elements such as those listed above. Computer-executable instructions may be compiled or interpreted from computer programs created using a variety of programming languages and/or technologies known to those skilled in the art, including, without limitation, and either alone or in combination, Java™, C, C++, Visual Basic, Java Script, Perl, etc. In general, a processor (e.g., a microprocessor) receives instructions, e.g., from a memory, a computer-readable medium, etc., and executes these instructions, thereby performing one or more processes, including one or more of the processes described herein. Such instructions and other data may be stored and transmitted using a variety of known computer-readable media.
A computer-readable medium includes any medium that participates in providing data (e.g., instructions), which may be read by a computer. Such a medium may take many forms, including, but not limited to, non-volatile media, and volatile media. Non-volatile media include, for example, optical or magnetic disks and other persistent memory. Volatile media include dynamic random access memory (DRAM), which typically constitutes a main memory. Common forms of computer-readable media include, for example, a floppy disk, a flexible disk, hard disk, magnetic tape, any other magnetic medium, a CD-ROM, DVD, any other optical medium, punch cards, paper tape, any other physical medium with patterns of holes, a RAM, a PROM, an EPROM, a FLASH-EEPROM, any other memory chip or cartridge, or any other medium from which a computer can read.
FIGS. 3-5 and the description thereof below present exemplary approaches to the functional details of the congestion control module 145. As illustrated, processes 300, 400 and 500 all operate concurrently. Concurrent operation may allow for the constant monitoring and detection of any failures in the servers 130 of the cluster 135. However, in another exemplary approach, the steps of the processes may be rearranged to operate sequentially. For example, after the initial set-up steps of process 300, process 400 may operate repeatedly for a period of time. Subsequently, process 400 may be paused while the cluster 135 is checked for server failures and process 500 updates the utilization factors.
FIG. 3 illustrates a flowchart of an exemplary process 300 for managing congestion control and server redundancy. The router 140 may include a computer-readable medium having stored instructions for carrying out certain operations described herein, including some or all of the operations described with respect to process 300. For example, some or all of such instructions may be included in the congestion control module 145. As described below, some steps of process 300 may include user input and interactions. However, it is to be understood that fully automated or other types of programmatic techniques may implement steps that include user input.
Process 300 begins in step 305 by receiving initial parameters. The initial parameters include at least the actual capacities 205 and redundancy factors for each of the servers 130. The parameters may be received from user input, e.g., via a command line interface, Graphical User Interface, etc. In another exemplary approach, the parameters may be provided in a configuration file, or the like. Accordingly, the parameters may be received in step 305 by opening the file for reading and extracting the relevant data.
Next, in step 310, an expected traffic load may be established. The expected traffic load may be used to alter or set the redundancy factors of the servers. Historical traffic loads for similar dates and times may be used to establish the expected traffic load. Moreover, the expected traffic load presets a baseline value to be used when initially setting the redundancy factors.
Next, in step 315, the actual capacity 205 of each server may be limited to respective available capacities 210 based on the redundancy factors. In generally, the sum of the available capacities 210, also referred to as the available cluster capacity 250, will correspond to the expected traffic load in order to ensure that all expected traffic will be able to be processed. Moreover, the limiting provides redundant capacity 215 which is reserved for times in which there is a failure of one of the servers 130 of the cluster 135. This initialization step sets baseline values for the available capacities 210. However, in the event of a server failure, the available capacities 210 may be altered by changing the respective redundancy factors.
Following, step 315, steps 320, 325, and 330 may operate concurrently as discussed above. Step 320 includes the steps and decisions of process 500 discussed below. Similarly, step 325 includes the steps and decisions of process 400, also discussed below. Because the utilization factor (uij) is based on the estimated current usage for the given time interval (j), the utilization factor will be calculated at each interval. The process will continue to schedule the data for the given time interval. At the conclusion of a time interval, the process must update the current usage based on the recent traffic load of each server.
In step 330, it is determined whether any of the servers 130 have failed. In one exemplary approach, the router 140 may attempt to contact the servers 130, e.g., by initiate a connection, transmitting a so-called ping (Internet Control Message Protocol echo), etc. In another exemplary approach, the servers 130 may be configured to send out a communication, sometimes referred to as a life beat, to the router 140. No matter the approach, the router 140 will continuously monitor the servers 130 for failures. For example, the lack of a response or the lack of a life beat may be indicative of a server failure.
In step 335, the redundancy factors are set to a failure state if a server failure is detected in step 340. As discussed above, the redundancy factors may be dynamically set to a high values, e.g., one, in order to allocate all of the redundant capacity 215 as available capacity 210.
In step 340, the redundancy factors are set to the initial parameters if a server failure is not detected. In most cases the redundancy factors will already be set to the initial parameters. However, if the functionally of a server 130 has just been restored following a failure, the redundancy factors may need to be changed from the failure state. As discussed above, the redundancy factors may be established such that the available cluster capacity 250 corresponds to the baseline or expected traffic load.
Next, in step 345, the parameters may be adjusted. For example, the redundancy factors may vary based on time and date to correspond with expected traffic loads and the service level that needs to be provided by the cluster 135. For example, if service must be guaranteed to a high degree for a certain time, the redundancy factors may be set to a low level to ensure there is redundant capacity 215 to accommodate any server failures. Accordingly, the parameters may be scheduled for particular times. However, in another exemplary approach, process 300 may be adaptive to current traffic conditions. For example, the parameters may automatically adjust in the face of changing traffic conditions.
If the parameters need to be adjusted, process 300 may return to step 305 to receive the new parameters. However, if no adjustment is required, process 300 may return to concurrent steps 320, 325, and 330.
FIG. 4 illustrates a flowchart of an exemplary process 400 for updating server utilization factors. The router 140 may include a computer-readable medium having stored instructions for carrying out certain operations described herein, including some or all of the operations described with respect to process 400. For example, some or all of such instructions may be included in the congestion control module 145. As discussed above, process 400 may be sub-process for process 300, e.g., in step 325.
Process 400 begins in step 405 by determining the available capacity 210 of each server. In one exemplary approach, the actual capacity 205 and the respective redundancy factor (ri) may be retrieved. For example, these values may be provided as initial parameters via user input, a configuration file, etc. The actual capacity 205 may be multiplied by the redundancy factor (1−ri) to determine the limited available capacity 210.
Next, in step 410, the current usage 220 of each server 130 may be estimated. Unlike the expected traffic load discussed above for setting baseline values for the available capacity, the estimated traffic load attempts to determine the current or actual traffic load. Because it may be too costly to constantly monitor the actual traffic load, e.g., amount of packet data traffic sent to a server 130, the router may break the monitoring into discrete time intervals. In one exemplary approach, the router 140 may monitor the amount of packet data sent to a server every 200 milliseconds. Accordingly, the router 140 only knows historical amounts of packet data traffic sent to a server 130. Moreover, the router may not know the actual amount of traffic sent to a sever during the instant time interval. However, because traffic volumes could potentially change dramatically even during a brief interval, an understanding of the current usage 220 of a server 130 is important for properly balancing the data processing load over the cluster 135. Moreover, the scheduling decision (discussed with respect to FIG. 5) is ideally based on an estimated current usage (Eλij) and not on historical usage.
In one exemplary approach, the current usage 220 may be based on a weighted moving average of the actual amounts of usage in prior intervals (λi,j−1), (λi,j−2), etc. Because the data processing may be based on sessions which exist and draw on server capacity for a period of time, it may be useful to base the current usage on more than just the actual usage of the most recent interval (λi,j−1). A weighting factor (0≦w≦1) may be selected to allocate the weight given to the most recent interval (λi,j l) and the second most recent interval (λi,j−2). For example, if it is unusual for a typical session to extend beyond a single interval, the weighting factor (w) may be set to a high value to give more weight the most recent period. Similarly, a lower value of (w) may be selected if it is likely that sessions draw on server capacity 210 for more than one time interval. Accordingly, the estimated current usage (Eλij) of each (i) server at each (j) time interval may be represented as Eλij=w·λi,j−1+(1−w)·λi,j−2. In other exemplary approaches, formulas that take even more historical values of the actual usage may be appropriate. As will be discussed with respect to FIG. 5, the actual usages for prior time intervals (λi,j−1) (λi,j−2), etc., may be stored during the packet data scheduling.
Next, in step 415, utilization factors for each server 130 may be calculated. The utilization factor may represent the estimated current usage (Eλij) as a ratio to the available capacity 210, where the available capacity is the actual capacity limited by the redundancy factor (ri). In one exemplary approach, the utilization factor (uij) may be expressed as
FIG. 5 illustrates a flowchart of an exemplary process 500 for scheduling packet data traffic to servers in the cluster. The router 140 may include a computer-readable medium having stored instructions for carrying out certain operations described herein, including some or all of the operations described with respect to process 500. For example, some or all of such instructions may be included in the congestion control module 145. As discussed above, process 500 may be sub-process for process 300, e.g., in step 320.
Process 500 begins in step 505 when incoming packet data traffic 155 is received by the router 140. As discussed above with respect to FIG. 2b, the router may classify, meter, mark, and shape the traffic as necessary. Subsequent to these preliminary steps, the scheduler 180, in coordination with the congestion control module 145, may proceed with the following steps to further process and direct the traffic to a particular server 130 of the cluster 135.
Next, in step 510, it is determined whether the received packet data belongs to an existing session. For example, the data may include a session identifier thereby associating the data with a particular session.
In step 515, the packet data will be scheduled to the server 130 that is already handling the session to which the data belongs. The session data from step 510 may be used to determine which server is handling a particular session. While not depicted, process 500 may also store associations between the servers 130 and sessions being processed thereby.
In step 520, following the determination in step 510 that the packet data does not belong to an existing session, the data will be scheduled to one of the servers 130 of the cluster 135. As discussed above with respect to process 400, utilization factors may be maintained and updated for each server. The utilization factors express the estimated usage of the server 130 with respect to the available capacity 210. The server with the highest utilization factor has the least amount of unused available capacity 210. To effectively balance the traffic load between the servers 130 of the cluster 135, the new traffic may be scheduled to the server 130 having the lowest utilization factor.
Following both steps 515 and 520, the record of the amount of packet data traffic sent to the server 130 may be updated in step 525. The router 140, or cluster manager, may keep historical records of the amounts of traffic sent to each server, e.g., a record for each of the last five time intervals. As discussed above, a time interval, e.g., 200 ms, may be defined for breaking down the calculation of the utilization factors. The historical records may be in the form of a circular list, or equivalent data structure, with an index value, e.g., 0-4, used to identify the record associated with the current time interval. Accordingly, the amount of the traffic scheduled to the server in either step 515 or step 520 will be added to the record associate with the server and the current time interval. While in another exemplary approach, the servers 130 could report back their actual current usage rates, monitoring estimated usage rates by the router 140 may eliminate time consuming bi-directional communication required for such reporting.
Accordingly, an exemplary congestion control method for session based network traffic has been described. Session handling servers 130 arranged in a cluster 135 may receive packet data traffic 155 from an upstream network device such as a router 140. The actual traffic handling capacity 205 of each server 130 may be limited by a redundancy factor (ri) to respective available capacities 210 in order to provide redundancy capacity 215. The redundancy factors (ri) may be altered dynamically in order to provide more or less redundancy in the cluster 135. For example, the redundancy factors (ri) may be decreased during times in which service availability is critically important. At other times, the redundancy factors may be increased to provide more available processing capacity 210, which may be useful in the event of a failure of one of the servers 130. To balance the traffic load across the cluster 135, the amount of data traffic sent to each server may be tracked and recorded for a number of historical time intervals. Some or all of the historical records may be used to estimate a current usage (Eλij). Newly received traffic that is not associated with an existing session may be scheduled to the server having the lowest utilization factor (uij), e.g., the ratio of the estimated usage to the available capacity.
With regard to the processes, systems, methods, heuristics, etc. described herein, it should be understood that, although the steps of such processes, etc. have been described as occurring according to a certain ordered sequence, such processes could be practiced with the described steps performed in an order other than the order described herein. It further should be understood that certain steps could be performed simultaneously, that other steps could be added, or that certain steps described herein could be omitted. In other words, the descriptions of processes herein are provided for the purpose of illustrating certain systems, and should in no way be construed so as to limit the claimed invention.
Accordingly, it is to be understood that the above description is intended to be illustrative and not restrictive. Many systems and applications other than the examples provided would be apparent upon reading the above description. The scope of the invention should be determined, not with reference to the above description, but should instead be determined with reference to the appended claims, along with the full scope of equivalents to which such claims are entitled. It is anticipated and intended that future developments will occur in the arts discussed herein, and that the disclosed systems and methods will be incorporated into such future systems. In sum, it should be understood that the disclosure is capable of modification and variation and is limited only by the following claims.
All terms used in the claims are intended to be given their broadest reasonable constructions and their ordinary meanings as understood by those skilled in the art unless an explicit indication to the contrary is made herein. In particular, use of the singular articles such as “a,” “the,” “said,” etc. should be read to recite one or more of the indicated elements unless a claim recites explicitly to the contrary. |
Geography of Massachusetts
Massachusetts, the 7th smallest state in the United States, resides in the New England region of the northeastern United States, and has an area of 10,555 square miles (27,340 km2). It is bordered on the north by New Hampshire and Vermont, on the west by New York, on the south by Connecticut and Rhode Island, and on the east by the Atlantic Ocean. It is the most populous New England state.
Massachusetts is called "the Bay State" because of several large bays, which distinctly shape its coast: Massachusetts Bay and Cape Cod Bay, to the east, and Buzzards Bay, to the south. A few cities and towns on the Massachusetts–Rhode Island border are adjacent to Narragansett Bay. At the southeastern corner of the state is a large, sandy, arm-shaped peninsula, Cape Cod. The islands Martha's Vineyard and Nantucket lie south of Cape Cod, across Nantucket Sound.
Boston is the largest city, at the inmost point of Massachusetts Bay, the mouth of the Charles River, the longest river entirely within Massachusetts. Most Bay Staters live in the Boston area, which cover most of eastern Massachusetts. Eastern Massachusetts is fairly densely populated and mostly suburban. Western Massachusetts is more rural and sparsely populated, especially in the Berkshires, the branch of the Appalachian Mountains that dominates the western quarter of the state. The most populous part of western Massachusetts is the Pioneer Valley, straddling the Connecticut River, which flows across Western Massachusetts from north to south.
Massachusetts has 351 cities and towns. Every part of the state is within an incorporated city or town, but many towns include large rural areas. The state's 14 counties have few government functions and serve as little more than judicial districts.
In Eastern Massachusetts, Boston is located at the innermost point of Massachusetts Bay, at the mouth of the Charles River. The Charles River is longest river located entirely within Massachusetts; however the Connecticut River in Western Massachusetts is the Commonwealth's longest river. Most of the population of the Boston metropolitan area (approximately 4.4 million) lives outside of the city proper. In general, Eastern Massachusetts, including and surrounding Boston, is densely populated. Boston's suburbs stretch as far west as the City of Worcester in Central Massachusetts.
Central Massachusetts encompasses Worcester County. It features the large city of Worcester, and the smaller cities of Fitchburg, Leominster, Gardner, and Southbridge. Central Massachusetts also includes many rural hill towns, forests, and small farms. The Quabbin Reservoir borders the western side of the county; it is the main water supply for Greater Boston.
West of the Central Massachusetts hill towns, the Pioneer Valley along the Connecticut River in Western Massachusetts features the Commonwealth's richest soil. The major city of Springfield sits beside the Connecticut River amidst a broad valley, a mere five miles (8 km) north of the Connecticut border and only 24 miles (39 km) from Connecticut's capital city, Hartford. The densely populated Springfield-Hartford region, called the Knowledge Corridor, is the second most populous region in New England (approximately 1.9 million.) As in Eastern Massachusetts, most residents live outside of the region's two principal cities, (i.e. Springfield and Hartford.) Other cities in the Massachusetts portion of the Knowledge Corridor include Chicopee, Agawam, West Springfield, Westfield, Holyoke, and the college towns of Northampton and Amherst.
West of the Knowledge Corridor is mountainous, including the hilltowns immediately to the west of the Valley. Further west rises a range of rolling, purple mountains known as the Berkshires. Near the New York border, the Taconic and Hoosac Ranges cross into Massachusetts; however, in general, the area is known as The Berkshires. The region was populated by aborigines until the 18th century when Scotch-Irish settlers arrived, after having found the fertile lowlands along the Connecticut River settled. On reaching the Berkshires, settlers found poor soil for farming, but discovered numerous fast-moving rivers for industry. Pittsfield and North Adams grew into small, prosperous cities. A number of smaller mill towns exist along the Westfield and Housatonic Rivers, interspersed among wealthy vacation resort towns.
The National Park Service administers a number of natural and historical sites in Massachusetts. Along with twelve national historic sites, areas, and corridors, the National Park Service also manages the Cape Cod National Seashore and the Boston Harbor Islands National Recreation Area.
In addition, the Massachusetts Department of Conservation and Recreation maintains a number of parks, trails, and beaches throughout the commonwealth.
Physical geography
Massachusetts extends from the mountains of the Appalachian system in the west to the sandy beaches and rocky shorelines of the Atlantic coast. The entire state was covered in ice during the Wisconsin glaciation, which shaped today’s landscape. Much of the state remains covered in glacial till and dotted with typical glacial features, such as kettle ponds, drumlins, eskers, and moraines. Apart from a few alluvial floodplains, soils tend to be rocky, acidic, and not very fertile.
Part of the state is uplands of resistant metamorphic rock that were scraped by Pleistocene glaciers that deposited moraines and outwash on a large, sandy, arm-shaped peninsula called Cape Cod and the islands Martha's Vineyard and Nantucket to the south of Cape Cod. Upland elevations increase dramatically in Western Massachusetts. These uplands are interrupted by the downfaulted southern Pioneer Valley along the Connecticut River and further west by the Housatonic Valley separating the Berkshire Hills from the Taconic Range along the western border with New York. The highest peak in the state is Mount Greylock at 3,491 feet (1,064 m) near the northwest corner.
Elevation and relief are greatest in the western part of the state and increase somewhat from south to north. The Taconic Mountains, part of the Appalachian system, run along the western border with New York, reaching 2,624 feet (800 meters) at Mount Everett in the state's southwest corner, and including the state’s highest point, Mount Greylock, at 3,491 feet (1,064 meters) in the northwest corner. The Housatonic-Hoosic valley separates the Taconics from The Berkshires, a broad belt of steeply rolling hills that are a southern extension of the Green Mountains of Vermont. They extend south to the border of Connecticut. Mount Greylock lies on the western edge of the Taconic Range, across the Hoosic River from the Hoosac Range to the east. The Hoosac Range connects the Green Mountains with the Berkshires.
Between the Berkshires and the rest of the state lies the Connecticut River Valley, known within Massachusetts as the Pioneer Valley. This ancient rift valley appeared in the Mesozoic Era when North and South America broke away from Europe and Africa. Dinosaur footprints near Mount Tom bear witness to that era, and series of basalt and sedimentary rock ridges (collectively known as the Metacomet Ridge) including Mount Toby, Mount Holyoke, Mount Tom, and others extending south to Long Island Sound and the valley's abrupt thousand-foot (300 meter) western escarpment illustrate the tectonic forces. More than a hundred million years later, as the Pleistocene epoch ended, receding glaciers left moraines that dammed the Connecticut River, creating Lake Hitchcock. Lacustrine silt deposits replaced soil scraped away by the glaciers, leaving behind deep, productive soil after the river breached the obstructing moraine and the lake disappeared.
East of this valley is an area of rolling uplands dotted with lakes and dissected by streams flowing into the Connecticut River in the west and into the Merrimack, Quinebaug, Blackstone, or Charles rivers, or into other shorter, coastal rivers in the east. Just to the east of the Pioneer Valley, hills rise steeply toward the divide between the Connecticut River basin and the river basins to the east. This divide runs through central Massachusetts, though the summit of Mount Wachusett, the highest point in the state east of the Connecticut River, rising to 2,006 feet (611 meters).
To the east of this divide, the elevation of the hilltops gradually decreases, and the landscape is more gently rolling. Within 30 miles (50 kilometers) of the coast, few hills exceed 300 feet (100 meters) in elevation. Near the coast, swamps, marshes, and ponds alternate with low hills. However, the Blue Hills, just south of Boston, rise above the surrounding landscape. The state probably takes its name from the Massachusett name for their highest point, Great Blue Hill, with an elevation of 635 feet (194 meters).
The Massachusetts coastline is deeply indented with bays, coves, and estuaries, separated by narrow promontories. Some of these form natural harbors that gave rise to the state’s historic ports, including Newburyport, Gloucester, Salem, Boston, and New Bedford. The state has a few small barrier islands, the largest of which is Plum Island. The state’s largest promontory is the Cape Cod peninsula. Its backbone is formed by glacial moraines, but much of its coastline has been shaped by the longshore drift of coastal sand, which forms many of its famous sandy beaches. To the south of Cape Cod, glacial moraines rise above the ocean surface to form the state’s largest islands: Martha’s Vineyard, Nantucket, the Elizabeth Islands, and Monomoy Island.
Massachusetts has a humid continental climate. Winters are cold, with average January temperatures below freezing nearly throughout the state, and summers are warm.
The hilly western interior of Massachusetts has the coldest winters. Stockbridge, in the Berkshires, has a January average temperature of 21.6°F (-5.8°C). Winters are more moderate along the eastern coast. Boston has the state's highest maximum January temperature—35.6°F (2 C), but this temperature is elevated by an urban heat island. The average January temperature in Hingham, also on the coast but 13 miles (21 km) southeast of Boston, is 27.5°F (-2.5°C). Summer temperatures are highest in the state's urban centers, due to the heat island effect. Boston's July temperature averages 81.7°F (27.6°C), and the July temperature in the central Massachusetts city of Worcester averages 79.2°F (26.2°C), which is measured at the Worcester Airport, at an elevation of over 1,000 feet (300m). By contrast, the coolest average summer temperatures occur in the Berkshires and on the state's offshore islands. The average temperature in August, the warmest month on Nantucket Island, is 68.7°F (20.4°C). The average in July in Stockbridge is 68.9°F (20.5°C). Both daily and seasonal variation in temperature are greatest in the western interior and lowest along the coast.
Precipitation is fairly evenly spread throughout the year in Massachusetts. Boston averages 43 in (1091 mm) of precipitation annually, with a maximum monthly average of 4.3 in (109.2 mm) in November and a minimum monthly average of 2.9 in (73.7 mm) in July. Springfield, in the Pioneer Valley, averages 45.8 in (1163.9 mm) of annual precipitation, with a 4.6 in (116.8 mm) maximum monthly average in June and a 2.7 in (68.6 mm) minimum monthly average in February. Interior Massachusetts tends to have a summer precipitation maximum due to convection in air masses heated over the interior, which gives rise to frequent thunderstorms. These occur less frequently over the coast, due to the relative lack of convection over the cooler ocean waters. On the other hand, cold, dry air masses over the interior of the state tend to suppress winter precipitation.
All of Massachusetts experiences substantial snowfall in a typical winter. Total annual snowfalls average 43.3 in (110.0 cm) in Boston and 69.1 in (175.5 cm) in Worcester. The ground is often covered with snow for weeks at a time in January and February.
Although Massachusetts has a humid climate, its climate is sunny compared to other humid climates at the same latitude. In Boston, the average percentage of possible sunshine for every month is at least 50%. In summer and early autumn, the average percentage of possible sunshine is greater than 60%, according to National Weather Service data. The hottest temperature recorded was 108 degrees Fahrenheit (42.8 degrees Celsius).
The primary biome of inland Massachusetts is temperate deciduous forest. Although much of the state had been cleared for agriculture, leaving only traces of old growth forest in isolated pockets, secondary growth has regenerated in many rural areas as farms have been abandoned. The areas most affected by human development include the Greater Boston area in the east, the smaller Springfield metropolitan area in the west, and the largely agricultural Pioneer Valley. Animals that have become locally extinct over the past few centuries include gray wolves, elk, wolverines, and mountain lions.
A number of species are doing well, despite, and in some cases because of the increased urbanization of the commonwealth. Peregrine falcons utilize office towers in larger cities as nesting areas, and the population of coyotes, whose diet may include garbage and roadkill, has been increasing in recent decades. White-tailed deer, raccoons, wild turkeys and eastern gray squirrels are also found throughout Massachusetts. In more rural areas in the western part of the state, larger mammals such as moose and black bears have returned, largely due to reforestation following the regional decline in agriculture.
Massachusetts is located along the Atlantic Flyway, a major route for migratory waterfowl along the Atlantic coast. Lakes in central Massachusetts provide habitat for the common loon, while a significant population of long-tailed ducks winter off Nantucket. Small offshore islands and beaches are home to roseate terns and are important breeding areas for the locally threatened piping plover. Protected areas such as the Monomoy National Wildlife Refuge provide critical breeding habitat for shorebirds and a variety of marine wildlife including a large population of gray seals.
Freshwater fish species in the commonwealth include bass, carp, catfish, and trout, while saltwater species such as Atlantic cod, haddock and American lobster populate offshore waters. Other marine species include Harbor seals, the endangered North Atlantic right whales, as well as humpback whales, fin whales, minke whales and Atlantic white-sided dolphins.
Most of Massachusetts is forested. Even suburban eastern Massachusetts is heavily wooded. Trees tend to grow around houses in this region, such that when one looks out over eastern Massachusetts from the top of a high hill, one sees a vista of treetops, punctuated only occasionally by a church steeple, smokestack, or radio tower.
According to U.S. government data , 46% of Massachusetts land is devoted to forest. Another 7% is rural parkland, which is also mainly forested. Urban and suburban development takes up 36% of the state’s land, but even this land, outside of the main urban centers, consists largely of houses on wooded properties. About 4% of the state’s land is cropland, and less than 1% is pasture. About 2% of the state’s land is marsh or other wetland. The remainder of the land is taken up with other uses, such as transportation.
Three ecoregions comprise the natural environment of Massachusetts. Atlantic coastal pine barrens occur on Cape Cod, Nantucket, and Martha's Vineyard. These are fire-prone temperate coniferous forests growing on the sandy soils of the coastal plain. The other two ecoregions are temperate broadleaf and mixed forests. Across most of the state, including eastern Massachusetts, south central Massachusetts, and the Connecticut River Valley, the Northeastern coastal forests are a mix of hardwood deciduous oak, maple, beech, and hickory and coniferous pine trees. In the Berkshires and north central Massachusetts, the more boreal New England-Acadian forests prevail. These consists mainly of coniferous spruce and hemlock, occasional pine, and deciduous birch trees. Roughly since the Civil War, farms have reverted to woodland. Lumbering activity has decreased in recent decades, so the more undisturbed forests have reclaimed some characteristics of old growth.
The forests (and wooded suburbs) are home to a variety of invertebrate and vertebrate animal species. The state has an abundance of white-tailed deer, and there have been concerns about deer overpopulation because many of the deer’s natural predators, such as wolves, have historically been hunted to extinction within Massachusetts. However, coyotes have been moving into Massachusetts to fill the ecological niche formerly occupied by wolves. Bears, wild turkey, and even moose have returned from northern refuges. In 1846 Thoreau traveled to Northern Maine to observe and write about moose, which he thought were well on the way to extinction. If he were alive today, he might find them almost within walking distance of Walden Pond.
Pollution, dams, and introduction of exotic species have decimated some native fish populations. Efforts to mitigate these problems and restore Atlantic salmon to the Connecticut River watershed have had very little success. The other widespread native salmonid, the brook trout, persists in cold upland streams, particularly above waterfalls and other barriers that exclude introduced brown and rainbow trout. American shad runs have retained at least a fraction of their former abundance, and smallmouth bass, sunfish, and pike populations are healthy enough to support angling.
Wetlands, including swamps and both salt- and fresh-water marshes, are important ecologically in Massachusetts. Many of the state’s fish and bird species inhabit wetland environments.
The state’s urban environments are partly wooded but also bear a heavy load of built structures and human environments that are not hospitable to many other species. At the same time, pollutants in waterways, mainly from urban sources, can be toxic to many species or may support algae and bacteria that lead to hypoxia and the death of aquatic animals. However, Greater Boston boasts extensive parklands, and efforts have been made in Massachusetts to reduce environmental pollution in both urban and rural parts of the state.
The Northeast megalopolis extends into Massachusetts. It occupies most of eastern Massachusetts starting at Worcester as well as the Springfield-Holyoke-Northampton urbanization that joins Connecticut's Hartford-New Haven urbanization.
According to the definitions of the U.S. Office of Management and Budget (OMB), all of Massachusetts falls within a metropolitan statistical area (MSA), except for the offshore islands of Martha’s Vineyard and Nantucket. According to 2005 Census estimates, 62% of the population of Massachusetts lives within the Boston MSA. Other Massachusetts metropolitan areas are the Worcester MSA (with 12% of the state's population), the Springfield MSA (11%), the Providence-Fall River-New Bedford MSA (9%), the Barnstable (Cape Cod) MSA (4%), and the Pittsfield MSA (2%).
In each of these metropolitan areas, population is concentrated in a number of densely populated cities and towns. In the Boston MSA, for example, the City of Boston and a cluster of densely populated inner suburbs within the Route 128 belt account for more than half of the population of the metropolitan area. The older cities of Lawrence, Lowell, and Brockton lie outside this urban core but are also densely populated.
However, population is growing fastest in the outer peripheries of the state's metropolitan areas, where new housing construction is adding dwelling units. While the state as a whole shows little population growth, or even a population decline in some years due to a net loss from migration, the belt of towns along Interstate 495, near the western edge of the Boston MSA, shows steady population growth.
The Springfield and Worcester MSAs include some very thinly populated rural areas. In the Berkshires and in the hills west of Worcester are a number of towns with population densities below 40 per square mile (compared with the state average of 810 per square mile).
Although the U.S. Census Bureau prepares population estimates for MSAs, these statistical units are defined by county borders. Because Massachusetts counties are relatively large and may contain several urban centers, MSAs are an imprecise way to describe the state’s urban clusters. For example, Lawrence, Lowell, and Brockton all have closer economic ties with neighboring towns than they do with one another. The Lowell region draws commuters from nearby New Hampshire who might not consider commuting all the way to Boston. Yet these areas are all part of the Boston MSA. Similarly, the cities of Leominster and Fitchburg form the core of a distinct urban cluster. Because they lie within Worcester County, however, they are considered part of the Worcester MSA.
Economic geography
A finer-grained statistical unit than the MSA is the New England City and Town Area, or NECTA. NECTAs take advantage of the administrative subdivision of the entire territory of Massachusetts and other New England states into towns and cities. (No part of Massachusetts is unincorporated county territory.) Each NECTA consists of a cluster of cities and towns defined by commuting patterns, which therefore correspond roughly to local labor markets. While the U.S. Census Bureau defines metropolitan areas by county boundaries, the U.S. Bureau of Labor Statistics (BLS) offers data on employment by NECTA.
By far the largest NECTA in Massachusetts is the Boston-Cambridge-Quincy (Greater Boston) NECTA, which covers eastern Massachusetts and extends into southern New Hampshire. This NECTA consists of a central Boston-Cambridge-Quincy NECTA Division, including the City of Boston and the surrounding cities and suburbs. The other satellite NECTA divisions in the Greater Boston NECTA are the Brockton-Bridgewater-Easton NECTA Division, the Framingham NECTA Division, the Haverhill-North Andover-Amesbury NECTA Division (extending well into southeastern New Hampshire), the Lawrence-Methuen-Salem NECTA Division (extending into southern New Hampshire), the Lowell-Billerica-Chelmsford (or Lowell) NECTA Division (extending into southern New Hampshire), the Lynn-Peabody-Salem NECTA Division, the Nashua NECTA Division (mainly in New Hampshire but including a few Massachusetts towns), and the Taunton-Norton-Raynham NECTA Division.
The other Massachusetts metropolitan NECTAs are the Barnstable Town NECTA (covering most of Cape Cod), the Leominster-Fitchburg-Gardner NECTA (in north central Massachusetts), the New Bedford NECTA (in southeastern Massachusetts), the Pittsfield NECTA (in far western Massachusetts), the Springfield NECTA (in the Pioneer Valley and extending into northern Connecticut), and the Worcester NECTA (in central Massachusetts, extending into northeastern Connecticut).
According to the BLS, total nonfarm employment in Massachusetts in 2005 was about 3.2 million. About half of these jobs were located in the Boston-Cambridge-Quincy NECTA Division, which lies entirely within Massachusetts, although this NECTA accounted for only about 43% of the state’s population, according to 2005 Census estimates. This indicates either a higher labor participation rate in central Greater Boston or a surplus of commuters traveling to work from other parts of Massachusetts or neighboring states. Clearly, Greater Boston dominates the employment and economy of Massachusetts.
The other major centers of employment in Massachusetts are the Springfield and Worcester NECTAs. The Springfield NECTA accounts for slightly more than 10% of the jobs in Massachusetts, while the Worcester NECTA accounts for slightly less than 10% of the state’s jobs. (Although both of these NECTAs extend into Connecticut, the towns that they include in Connecticut account for only a small portion of their population and, probably, of their employment).
In every Massachusetts NECTA, service-sector jobs far outnumber goods-producing (natural resources, construction, and manufacturing) jobs. Beyond this generalization, there are some differences in the employment and economic structures of the state’s NECTAs and NECTA divisions.
In the far southeastern corner of Massachusetts, the Barnstable Town NECTA, nearly coterminous with the summer resort region of Cape Cod, has an atypical employment structure. It has the lowest share of employment in goods-producing jobs, which account for only 9.5% of its employment. Most of these jobs are in the construction sector. Manufacturing jobs account for only 3.3% of employment, compared with 9.6% for the state as a whole. On the other hand, the Cape Cod NECTA has the state’s highest percentages of employment in retail trade (17.9%, versus 11.1% for the state) and in leisure and hospitality (16.9%, versus 9.1% for the state). These numbers reflect the continuing importance to Cape Cod of summer tourism.
The central Boston-Cambridge-Quincy division of the larger NECTA with the same name also has a relatively low percentage (6.7%) of manufacturing employment. Although this division accounts for about half of the state’s total employment, it has only about a third of the state’s manufacturing jobs. Its largest manufacturing subsector is the production of computers and electronic products (28% of the division’s manufacturing jobs). This subsector is centered not in Boston’s urban core, but in the suburbs to the north and west, along Route 128. The economy of central Greater Boston is even more biased toward service provision than that of the rest of the state.
The particular economic strength of central Greater Boston is knowledge-intensive activities. It accounts for 62.2% of the state’s information sector jobs, and 66.0% of the jobs in the software-publishing subsector. Central Greater Boston has 68.8% of the state’s financial sector jobs, and 92.5% of the jobs in the investment subsector. It has 69.3% of the state’s jobs in management and technical consulting. Greater Boston is noted nationwide for its prestigious institutions of higher education, such as Harvard University and MIT, and the region is home to 77.8% of the state’s higher-education employment. Together, the knowledge-intensive information, financial, professional and business services, and education sectors account for 36.6% of the jobs in central Greater Boston, compared with 28.8% of the jobs in Massachusetts as a whole and 23.2% for the United States as a whole.
The satellite NECTA divisions that lie on the periphery of the Greater Boston NECTA all have higher percentages of employment in manufacturing than central Greater Boston or than Massachusetts as a whole. Many of these satellite NECTA divisions are centered on historic manufacturing cities, such as Haverhill, Lawrence, Lowell, Lynn, and Brockton. The BLS breaks down manufacturing employment only for the Framingham and Lowell NECTA divisions, to the west and northwest of Boston, respectively. In both of these divisions, computer and electronics manufacturing accounts for well over half of manufacturing employment. Except for Lowell, these satellite NECTA divisions also have higher shares of employment in retail trade than central Greater Boston or Massachusetts as a whole. These divisions, located along the major highways radiating from Boston, are particularly rich in shopping centers and wholesalers. The Lowell and Framingham divisions have even higher shares of employment in the information sector than central Greater Boston. This reflects the strength of these regions in the software publishing and telecommunications subsectors. On the other hand, these satellite divisions have lower shares of employment in financial services and in health and education services than the state average, reflecting the regional dominance of central Greater Boston in these areas. The Framingham division, however, has the state’s highest percentage of jobs in professional and business services (18.5% of employment versus 14.4% statewide), reflecting that region’s strength in technology.
The New Bedford NECTA has the state’s second-highest percentage (16.6%) of manufacturing employment. It has the state’s lowest percentages of employment in the financial sector (3.1%) and in professional and business services (6.25%).
The Leominster-Fitchburg-Gardner NECTA has the state’s highest percentage (17.8%) of manufacturing employment. It has by far the state’s lowest percentage of employment (1.0%) in the information sector and the second-lowest rate of employment in professional and business services (6.73%). On the other hand, this NECTA has the state’s highest percentage of employment (16.4%) in government.
The Worcester NECTA has a relatively high percentage (12.0%) of employment in manufacturing. Next to the Barnstable Town NECTA, it has a high percentage (14.9%) of employment in the healthcare sector. It has the lowest percentage of employment (8.7%) in the leisure and hospitality sector, reflecting the relative underdevelopment of its tourism industry.
The Springfield NECTA also has relatively high (12.9%) manufacturing employment. It has the state’s largest percentage of employment in the transportation and utilities subsector (4.5%, versus 2.6% for the state as a whole). It has the second-highest percentage (16.3%) of jobs in government.
Despite the small size of the Pittsfield NECTA, its employment by sector is similar to that of Massachusetts as a whole for most sectors. However, it has the state’s highest percentage of employment (20.4%) in the education and healthcare sector. It also has the second-highest share of employment (13.2%) in the leisure and hospitality sector. This reflects the importance of tourism in the Berkshires to the region’s economy.
See also
- List of mountains in Massachusetts
- List of Massachusetts rivers
- Northern boundary of Massachusetts
- "Population, Housing Units, Area, and Density (geographically ranked by total population): 2000". United States Census Bureau. Retrieved 2010-05-30.
- "Charles River Watershed". Office of Energy and Environmental Affairs. Retrieved 2010-05-23.
- The North Quabbin Woods: www.northquabbinwoods.org
- PDF (390 KB) (map; see text on map). Secretary of the Commonwealth of Massachusetts. Retrieved January 14, 2007.
- "Massachusetts". National Park Service. Retrieved 2010-05-26.
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- "Getting Wet!". Massachusetts Department of Conservation and Recreation. Retrieved 2010-05-26.
- "Elevations and Distances in the United States". U.S Geological Survey. 29 April 2005. Retrieved November 6, 2006.
- Forest Physiography: Physiography of the United States and Principles of Soils in Relation to Forestry, Isaiah Bowman (New York: Wiley and Sons, 1911): p. 681.
- "A Short Introduction to Terrestrial Biomes". www.nearctica.com. Retrieved 2009-10-17.
- Stocker, Carol. Old growth, grand specimens drive big-tree hunters The Boston Globe. Novemberttp://www.umass.edu/ruralmass/currentresearch.html
- "Massachusetts Forests". MassWoods Forest Conservation Program — The University of Massachusetts. Retrieved 2009-03-19.
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- "State Mammal List". Massachusetts Division of Fisheries & Wildlife. Retrieved 2009-10-17.
- "Peregrine Falcon". Massachusetts Division of Fisheries & Wildlife. Retrieved 2010-05-26.
- "Eastern Coyote in Massachusetts". Massachusetts Division of Fisheries & Wildlife. Retrieved 2010-05-26.
- "Wild Turkey in Massachusetts". Massachusetts Division of Fisheries & Wildlife. Retrieved 2010-05-26.
- "Moose in Massachusetts". Massachusetts Division of Fisheries & Wildlife. Retrieved 2010-05-26.
- "Black Bears in Massachusetts". Massachusetts Division of Fisheries & Wildlife. Retrieved 2010-05-26.
- "Atlantic Flyway". University of Nebraska. Retrieved 2010-05-22.
- "Common Loon". Massachusetts Division of Fisheries & Wildlife. Retrieved 2010-05-28.
- "Telemetry Research:Long-Tailed Ducks". Mass Audubon. Retrieved 2010-05-28.
- "Roseate Tern". Massachusetts Division of Fisheries & Wildlife. Retrieved 2010-05-28.
- "Coastal Waterbird Program". Mass Audubon. Retrieved 2010-05-28.
- "Monomoy National Wildlife Refuge - Wildlife and Habitat". United States Fish and Wildlife Service. Retrieved 2010-05-26.
- "Best Bets for Fishing". Massachusetts Division of Wildlife & Fisheries. Retrieved 2010-05-30.
- "Species Profiles". Massachusetts Division of Marine Fisheries. Retrieved 2010-05-30.
- Olson, D. M, E. Dinerstein, et al (2001). "Terrestrial Ecoregions of the World: A New Map of Life on Earth". BioScience 51 (11): 933–938. doi:10.1641/0006-3568(2001)051[0933:TEOTWA]2.0.CO;2.
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The term ‘homosexuality’ was coined in the late 19th century by a German psychologist, Karoly Maria Benkert. Although the term is new, discussions about sexuality in general, and same-sex attraction in particular, have occasioned philosophical discussion ranging from Plato's Symposium to contemporary queer theory. Since the history of cultural understandings of same-sex attraction is relevant to the philosophical issues raised by those understandings, it is necessary to review briefly some of the social history of homosexuality. Arising out of this history, at least in the West, is the idea of natural law and some interpretations of that law as forbidding homosexual sex. References to natural law still play an important role in contemporary debates about homosexuality in religion, politics, and even courtrooms. Finally, perhaps the most significant recent social change involving homosexuality is the emergence of the gay liberation movement in the West. In philosophical circles this movement is, in part, represented through a rather diverse group of thinkers who are grouped under the label of queer theory. A central issue raised by queer theory, which will be discussed below, is whether homosexuality, and hence also heterosexuality and bisexuality, is socially constructed or purely driven by biological forces.
- 1. History
- 2. Historiographical Debates
- 3. Natural Law
- 4. Queer Theory and the Social Construction of Sexuality
- 5. Conclusion
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- Related Entries
As has been frequently noted, the ancient Greeks did not have terms or concepts that correspond to the contemporary dichotomy of ‘heterosexual’ and ‘homosexual’. There is a wealth of material from ancient Greece pertinent to issues of sexuality, ranging from dialogues of Plato, such as the Symposium, to plays by Aristophanes, and Greek artwork and vases. What follows is a brief description of ancient Greek attitudes, but it is important to recognize that there was regional variation. For example, in parts of Ionia there were general strictures against same-sex eros, while in Elis and Boiotia (e.g., Thebes), it was approved of and even celebrated (cf. Dover, 1989; Halperin, 1990).
Probably the most frequent assumption of sexual orientation is that persons can respond erotically to beauty in either sex. Diogenes Laeurtius, for example, wrote of Alcibiades, the Athenian general and politician of the 5th century B.C., “in his adolescence he drew away the husbands from their wives, and as a young man the wives from their husbands.” (Quoted in Greenberg, 1988, 144) Some persons were noted for their exclusive interests in persons of one gender. For example, Alexander the Great and the founder of Stoicism, Zeno of Citium, were known for their exclusive interest in boys and other men. Such persons, however, are generally portrayed as the exception. Furthermore, the issue of what gender one is attracted to is seen as an issue of taste or preference, rather than as a moral issue. A character in Plutarch's Erotikos (Dialogue on Love) argues that “the noble lover of beauty engages in love wherever he sees excellence and splendid natural endowment without regard for any difference in physiological detail.” (Ibid., 146) Gender just becomes irrelevant “detail” and instead the excellence in character and beauty is what is most important.
Even though the gender that one was erotically attracted to (at any specific time, given the assumption that persons will likely be attracted to persons of both sexes) was not important, other issues were salient, such as whether one exercised moderation. Status concerns were also of the highest importance. Given that only free men had full status, women and male slaves were not problematic sexual partners. Sex between freemen, however, was problematic for status. The central distinction in ancient Greek sexual relations was between taking an active or insertive role, versus a passive or penetrated one. The passive role was acceptable only for inferiors, such as women, slaves, or male youths who were not yet citizens. Hence the cultural ideal of a same-sex relationship was between an older man, probably in his 20's or 30's, known as the erastes, and a boy whose beard had not yet begun to grow, the eromenos or paidika. In this relationship there was courtship ritual, involving gifts (such as a rooster), and other norms. The erastes had to show that he had nobler interests in the boy, rather than a purely sexual concern. The boy was not to submit too easily, and if pursued by more than one man, was to show discretion and pick the more noble one. There is also evidence that penetration was often avoided by having the erastes face his beloved and place his penis between the thighs of the eromenos, which is known as intercrural sex. The relationship was to be temporary and should end upon the boy reaching adulthood (Dover, 1989). To continue in a submissive role even while one should be an equal citizen was considered troubling, although there certainly were many adult male same-sex relationships that were noted and not strongly stigmatized. While the passive role was thus seen as problematic, to be attracted to men was often taken as a sign of masculinity. Greek gods, such as Zeus, had stories of same-sex exploits attributed to them, as did other key figures in Greek myth and literature, such as Achilles and Hercules. Plato, in the Symposium, argues for an army to be comprised of same-sex lovers. Thebes did form such a regiment, the Sacred Band of Thebes, formed of 500 soldiers. They were renowned in the ancient world for their valor in battle.
Ancient Rome had many parallels in its understanding of same-sex attraction, and sexual issues more generally, to ancient Greece. This is especially true under the Republic. Yet under the Empire, Roman society slowly became more negative in its views towards sexuality, probably due to social and economic turmoil, even before Christianity became influential.
Exactly what attitude the New Testament has towards sexuality in general, and same-sex attraction in particular, is a matter of sharp debate. John Boswell argues, in his fascinating Christianity, Social Tolerance, and Homosexuality, that many passages taken today as condemnations of homosexuality are more concerned with prostitution, or where same-sex acts are described as “unnatural” the meaning is more akin to ‘out of the ordinary’ rather than as immoral (Boswell, 1980, ch.4; see also Boswell, 1994). Yet others have criticized, sometimes persuasively, Boswell's scholarship (see Greenberg, 1988, ch.5). What is clear, however, is that while condemnation of same-sex attraction is marginal to the Gospels and only an intermittent focus in the rest of the New Testament, early Christian church fathers were much more outspoken. In their writings there is a horror at any sort of sex, but in a few generations these views eased, in part due no doubt to practical concerns of recruiting converts. By the fourth and fifth centuries the mainstream Christian view allowed for procreative sex.
This viewpoint, that procreative sex within marriage is allowed, while every other expression of sexuality is sinful, can be found, for example, in St. Augustine. This understanding leads to a concern with the gender of one's partner that is not found in previous Greek or Roman views, and it clearly forbids homosexual acts. Soon this attitude, especially towards homosexual sex, came to be reflected in Roman Law. In Justinian's Code, promulgated in 529, persons who engaged in homosexual sex were to be executed, although those who were repentant could be spared. Historians agree that the late Roman Empire saw a rise in intolerance towards sexuality, although there were again important regional variations.
With the decline of the Roman Empire, and its replacement by various barbarian kingdoms, a general tolerance (with the sole exception of Visigothic Spain) of homosexual acts prevailed. As one prominent scholar puts it, “European secular law contained few measures against homosexuality until the middle of the thirteenth century.” (Greenberg, 1988, 260) Even while some Christian theologians continued to denounce nonprocreative sexuality, including same-sex acts, a genre of homophilic literature, especially among the clergy, developed in the eleventh and twelfth centuries (Boswell, 1980, chapters 8 and 9).
The latter part of the twelfth through the fourteenth centuries, however, saw a sharp rise in intolerance towards homosexual sex, alongside persecution of Jews, Muslims, heretics, and others. While the causes of this are somewhat unclear, it is likely that increased class conflict alongside the Gregorian reform movement in the Catholic Church were two important factors. The Church itself started to appeal to a conception of “nature” as the standard of morality, and drew it in such a way so as to forbid homosexual sex (as well as extramarital sex, nonprocreative sex within marriage, and often masturbation). For example, the first ecumenical council to condemn homosexual sex, Lateran III of 1179, stated that “Whoever shall be found to have committed that incontinence which is against nature” shall be punished, the severity of which depended upon whether the transgressor was a cleric or layperson (quoted in Boswell, 1980, 277). This appeal to natural law (discussed below) became very influential in the Western tradition. An important point to note, however, is that the key category here is the ‘sodomite,’ which differs from the contemporary idea of ‘homosexual’. A sodomite was understood as act-defined, rather than as a type of person. Someone who had desires to engage in sodomy, yet did not act upon them, was not a sodomite. Also, persons who engaged in heterosexual sodomy were also sodomites. There are reports of persons being burned to death or beheaded for sodomy with a spouse (Greenberg, 1988, 277). Finally, a person who had engaged in sodomy, yet who had repented of his sin and vowed to never do it again, was no longer a sodomite. The gender of one's partner is again not of decisive importance, although some medieval theologians single out same-sex sodomy as the worst type of sexual crime.
For the next several centuries in Europe, the laws against homosexual sex were severe in their penalties. Enforcement, however, was episodic. In some regions, decades would pass without any prosecutions. Yet the Dutch, in the 1730's, mounted a harsh anti-sodomy campaign (alongside an anti-Gypsy pogrom), even using torture to obtain confessions. As many as one hundred men and boys were executed and denied burial (Greenberg, 1988, 313-4). Also, the degree to which sodomy and same-sex attraction were accepted varied by class, with the middle class taking the narrowest view, while the aristocracy and nobility often accepted public expressions of alternative sexualities. At times, even with the risk of severe punishment, same-sex oriented subcultures would flourish in cities, sometimes only to be suppressed by the authorities. In the 19th century there was a significant reduction in the legal penalties for sodomy. The Napoleonic code decriminalized sodomy, and with Napoleon's conquests that Code spread. Furthermore, in many countries where homosexual sex remained a crime, the general movement at this time away from the death penalty usually meant that sodomy was removed from the list of capital offenses.
In the 18th and 19th centuries an overtly theological framework no longer dominated the discourse about same-sex attraction. Instead, secular arguments and interpretations became increasingly common. Probably the most important secular domain for discussions of homosexuality was in medicine, including psychology. This discourse, in turn, linked up with considerations about the state and its need for a growing population, good soldiers, and intact families marked by clearly defined gender roles. Doctors were called in by courts to examine sex crime defendants (Foucault, 1980; Greenberg, 1988). At the same time, the dramatic increase in school attendance rates and the average length of time spent in school, reduced transgenerational contact, and hence also the frequency of transgenerational sex. Same-sex relations between persons of roughly the same age became the norm.
Clearly the rise in the prestige of medicine resulted in part from the increasing ability of science to account for natural phenomena on the basis of mechanistic causation. The application of this viewpoint to humans led to accounts of sexuality as innate or biologically driven. The voluntarism of the medieval understanding of sodomy, that sodomites chose sin, gave way to the modern notion of homosexuality as a deep, unchosen characteristic of persons, regardless of whether they act upon that orientation. The idea of a ‘latent sodomite’ would not have made sense, yet under this new view it does make sense to speak of a person as a ‘latent homosexual.’ Instead of specific acts defining a person, as in the medieval view, an entire physical and mental makeup, usually portrayed as somehow defective or pathological, is ascribed to the modern category of ‘homosexual.’ Although there are historical precursors to these ideas (e.g., Aristotle gave a physiological explanation of passive homosexuality), medicine gave them greater public exposure and credibility (Greenberg, 1988, ch.15). The effects of these ideas cut in conflicting ways. Since homosexuality is, by this view, not chosen, it makes less sense to criminalize it. Persons are not choosing evil acts. Yet persons may be expressing a diseased or pathological mental state, and hence medical intervention for a cure is appropriate. Hence doctors, especially psychiatrists, campaigned for the repeal or reduction of criminal penalties for consensual homosexual sodomy, yet intervened to “rehabilitate” homosexuals. They also sought to develop techniques to prevent children from becoming homosexual, for example by arguing that childhood masturbation caused homosexuality, hence it must be closely guarded against.
In the 20th century sexual roles were redefined once again. For a variety of reasons, premarital intercourse slowly became more common and eventually acceptable. With the decline of prohibitions against sex for the sake of pleasure even outside of marriage, it became more difficult to argue against gay sex. These trends were especially strong in the 1960's, and it was in this context that the gay liberation movement took off. Although gay and lesbian rights groups had been around for decades, the low-key approach of the Mattachine Society (named after a medieval secret society) and the Daughters of Bilitis had not gained much ground. This changed in the early morning hours of June 28, 1969, when the patrons of the Stonewall Inn, a gay bar in Greenwich Village, rioted after a police raid. In the aftermath of that event, gay and lesbian groups began to organize around the country. Gay Democratic clubs were created in every major city, and one fourth of all college campuses had gay and lesbian groups (Shilts, 1993, ch.28). Large gay urban communities in cities from coast to coast became the norm. The American Psychiatric Association removed homosexuality from its official listing of mental disorders. The increased visibility of gays and lesbians has become a permanent feature of American life despite the two critical setbacks of the AIDS epidemic and an anti-gay backlash (see Berman, 1993, for a good survey). The post-Stonewall era has also seen marked changes in Western Europe, where the repeal of anti-sodomy laws and legal equality for gays and lesbians has become common.
Broader currents in society have influenced the ways in which scholars and activists have approached research into sexuality and same-sex attraction. Some early 20th century researchers and equality advocates, seeking to vindicate same-sex relations in societies that disparaged and criminalized it, put forward lists of famous historical figures attracted to persons of the same sex. Such lists implied a common historical entity underlying sexual attraction, whether one called it ‘inversion’ or ‘homosexuality.’ This approach (or perhaps closely related family of approaches) is commonly called essentialism. Historians and researchers sympathetic to the gay liberation movement of the late 1960s and 1970s produced a number of books that implicitly relied on an essentialist approach. In the 1970s and 1980s John Boswell raised it to a new level of methodological and historical sophistication, although his position shifted over time to one of virtual agnosticism between essentialists and their critics. Crompton’s work (2003) is a notable contemporary example of an essentialist methodology.
Essentialists claim that categories of sexual attraction are observed rather than created. For example, while ancient Greece did not have terms that correspond to the heterosexual/homosexual division, persons did note men who were only attracted to person of a specific sex. Through history and across cultures there are consistent features, albeit with meaningful variety over time and space, in sexual attraction to the point that it makes sense of speak of specific sexual orientations. According to this view, homosexuality is a specific, natural kind rather than a cultural or historical product. Essentialists allow that there are cultural differences in how homosexuality is expressed and interpreted, but they emphasize that this does not prevent it from being a universal category of human sexual expression.
In contrast, in the 1970s and since a number of researchers, often influenced by Mary McIntosh or Michel Foucault, argued that class relations, the human sciences, and other historically constructed forces create sexual categories and the personal identities associated with them. For advocates of this view, such as David Halperin, how sex is organized in a given cultural and historical setting is irreducibly particular (Halperin, 2002). The emphasis on the social creation of sexual experience and expression led to the labeling of the viewpoint as social constructionism, although more recently several of its proponents have preferred the term ‘historicism.’ Thus homosexuality, as a specific sexual construction, is best understood as a solely modern, Western concept and role. Prior to the development of this construction, persons were not really ‘homosexual’ even when they were only attracted to persons of the same sex. The differences between, say, ancient Greece, with its emphasis on pederasty, role in the sex act, and social status, and the contemporary Western role of ‘gay’ or ‘homosexual’ are simply too great to collapse into one category.
In a manner closely related to the claims of queer theory, discussed below, social constructionists argue that specific social constructs produce sexual ways of being. There is no given mode of sexuality that is independent of culture; even the concept and experience of sexual orientation itself are products of history. For advocates of this view, the range of historical sexual diversity, and the fluidity of human possibility, is simply too varied to be adequately captured by any specific conceptual scheme.
There is a significant political dimension to this seemingly abstract historiographical debate. Social constructionists argue that essentialism is the weaker position politically for at least two reasons. First, by accepting a basic heterosexual/homosexual organizing dichotomy, essentialism wrongly concedes that heterosexuality is the norm and that homosexuality is, strictly speaking, abnormal and the basis for a permanent minority. Second, social constructionists argue that an important goal of historical investigations should be to put into question contemporary organizing schemas about sexuality. The acceptance of the contemporary heterosexual/homosexual dichotomy is conservative, perhaps even reactionary, and forecloses the exploration of new possibilities. (There are related queer theory criticisms of the essentialist position, discussed below.) In contrast, essentialists argue that a historicist approach forecloses the very possibility of a ‘gay history.’ Instead, the field of investigation becomes other social forces and how they ‘produce’ a distinct form or forms of sexuality. Only an essentialist approach can maintain the project of gay history, and minority histories in general, as a force for liberation.
Today natural law theory offers the most common intellectual defense for differential treatment of gays and lesbians, and as such it merits attention. The development of natural law is a long and very complicated story, but a reasonable place to begin is with the dialogues of Plato, for this is where some of the central ideas are first articulated, and, significantly enough, are immediately applied to the sexual domain. For the Sophists, the human world is a realm of convention and change, rather than of unchanging moral truth. Plato, in contrast, argued that unchanging truths underpin the flux of the material world. Reality, including eternal moral truths, is a matter of phusis. Even though there is clearly a great degree of variety in conventions from one city to another (something ancient Greeks became increasingly aware of), there is still an unwritten standard, or law, that humans should live under.
In the Laws, Plato applies the idea of a fixed, natural law to sex, and takes a much harsher line than he does in the Symposium or the Phraedrus. In Book One he writes about how opposite-sex sex acts cause pleasure by nature, while same-sex sexuality is “unnatural” (636c). In Book Eight, the Athenian speaker considers how to have legislation banning homosexual acts, masturbation, and illegitimate procreative sex widely accepted. He then states that this law is according to nature (838-839d). Probably the best way of understanding Plato's discussion here is in the context of his overall concerns with the appetitive part of the soul and how best to control it. Plato clearly sees same-sex passions as especially strong, and hence particularly problematic, although in the Symposium that erotic attraction could be the catalyst for a life of philosophy, rather than base sensuality (Cf. Dover, 1989, 153-170; Nussbaum, 1999, esp. chapter 12).
Other figures played important roles in the development of natural law theory. Aristotle, with his emphasis upon reason as the distinctive human function, and the Stoics, with their emphasis upon human beings as a part of the natural order of the cosmos, both helped to shape the natural law perspective which says that “True law is right reason in agreement with nature,” as Cicero put it. Aristotle, in his approach, did allow for change to occur according to nature, and therefore the way that natural law is embodied could itself change with time, which was an idea Aquinas later incorporated into his own natural law theory. Aristotle did not write extensively about sexual issues, since he was less concerned with the appetites than Plato. Probably the best reconstruction of his views places him in mainstream Greek society as outlined above; the main issue is that of active versus a passive role, with only the latter problematic for those who either are or will become citizens. Zeno, the founder of Stoicism, was, according to his contemporaries, only attracted to men, and his thought had no prohibitions against same-sex sexuality. In contrast, Cicero, a later Stoic, was dismissive about sexuality in general, with some harsher remarks towards same-sex pursuits (Cicero, 1966, 407-415).
The most influential formulation of natural law theory was made by Thomas Aquinas in the thirteenth century. Integrating an Aristotelian approach with Christian theology, Aquinas emphasized the centrality of certain human goods, including marriage and procreation. While Aquinas did not write much about same-sex sexual relations, he did write at length about various sex acts as sins. For Aquinas, sexuality that was within the bounds of marriage and which helped to further what he saw as the distinctive goods of marriage, mainly love, companionship, and legitimate offspring, was permissible, and even good. Aquinas did not argue that procreation was a necessary part of moral or just sex; married couples could enjoy sex without the motive of having children, and sex in marriages where one or both partners is sterile (perhaps because the woman is postmenopausal) is also potentially just (given a motive of expressing love). So far Aquinas' view actually need not rule out homosexual sex. For example, a Thomist could embrace same-sex marriage, and then apply the same reasoning, simply seeing the couple as a reproductively sterile, yet still fully loving and companionate union.
Aquinas, in a significant move, adds a requirement that for any given sex act to be moral it must be of a generative kind. The only way that this can be achieved is via vaginal intercourse. That is, since only the emission of semen in a vagina can result in natural reproduction, only sex acts of that type are generative, even if a given sex act does not lead to reproduction, and even if it is impossible due to infertility. The consequence of this addition is to rule out the possibility, of course, that homosexual sex could ever be moral (even if done within a loving marriage), in addition to forbidding any non-vaginal sex for opposite-sex married couples. What is the justification for this important addition? This question is made all the more pressing in that Aquinas does allow that how broad moral rules apply to individuals may vary considerably, since the nature of persons also varies to some extent. That is, since Aquinas allows that individual natures vary, one could simply argue that one is, by nature, emotionally and physically attracted to persons of one's own gender, and hence to pursue same-sex relationships is ‘natural’ (Sullivan, 1995). Unfortunately, Aquinas does not spell out a justification for this generative requirement.
More recent natural law theorists, however, have tried a couple different lines of defense for Aquinas' ‘generative type’ requirement. The first is that sex acts that involve either homosexuality, heterosexual sodomy, or which use contraception, frustrate the purpose of the sex organs, which is reproductive. This argument, often called the ‘perverted faculty argument’, is perhaps implicit in Aquinas. It has, however, come in for sharp attack (see Weitham, 1997), and the best recent defenders of a Thomistic natural law approach are attempting to move beyond it (e.g., George, 1999, dismisses the argument). If their arguments fail, of course, they must allow that some homosexual sex acts are morally permissible (even positively good), although they would still have resources with which to argue against casual gay (and straight) sex.
Although the specifics of the second sort of argument offered by various contemporary natural law theorists vary, the common elements are strong (Finnis, 1994; George, 1999). As Thomists, their argument rests largely upon an account of human goods. The two most important for the argument against homosexual sex (though not against homosexuality as an orientation which is not acted upon, and hence in this they follow official Catholic doctrine; see George, 1999, ch.15) are personal integration and marriage. Personal integration, in this view, is the idea that humans, as agents, need to have integration between their intentions as agents and their embodied selves. Thus, to use one's or another's body as a mere means to one's own pleasure, as they argue happens with masturbation, causes ‘dis-integration’ of the self. That is, one's intention then is just to use a body (one's own or another's) as a mere means to the end of pleasure, and this detracts from personal integration. Yet one could easily reply that two persons of the same sex engaging in sexual union does not necessarily imply any sort of ‘use’ of the other as a mere means to one's own pleasure. Hence, natural law theorists respond that sexual union in the context of the realization of marriage as an important human good is the only permissible expression of sexuality. Yet this argument requires drawing how marriage is an important good in a very particular way, since it puts procreation at the center of marriage as its “natural fulfillment” (George, 1999, 168). Natural law theorists, if they want to support their objection to homosexual sex, have to emphasize procreation. If, for example, they were to place love and mutual support for human flourishing at the center, it is clear that many same-sex couples would meet this standard. Hence their sexual acts would be morally just.
There are, however, several objections that are made against this account of marriage as a central human good. One is that by placing procreation as the ‘natural fulfillment’ of marriage, sterile marriages are thereby denigrated. Sex in an opposite-sex marriage where the partners know that one or both of them are sterile is not done for procreation. Yet surely it is not wrong. Why, then, is homosexual sex in the same context (a long-term companionate union) wrong (Macedo, 1995)? The natural law rejoinder is that while vaginal intercourse is a potentially procreative sex act, considered in itself (though admitting the possibility that it may be impossible for a particular couple), oral and anal sex acts are never potentially procreative, whether heterosexual or homosexual (George, 1999). But is this biological distinction also morally relevant, and in the manner that natural law theorists assume? Natural law theorists, in their discussions of these issues, seem to waver. On the one hand, they want to defend an ideal of marriage as a loving union wherein two persons are committed to their mutual flourishing, and where sex is a complement to that ideal. Yet that opens the possibility of permissible gay sex, or heterosexual sodomy, both of which they want to oppose. So they then defend an account of sexuality which seems crudely reductive, emphasizing procreation to the point where literally a male orgasm anywhere except in the vagina of one's loving spouse is impermissible. Then, when accused of being reductive, they move back to the broader ideal of marriage.
Natural law theory, at present, has made significant concessions to mainstream liberal thought. In contrast certainly to its medieval formulation, most contemporary natural law theorists argue for limited governmental power, and do not believe that the state has an interest in attempting to prevent all moral wrongdoing. Still, they do argue against homosexuality, and against legal protections for gays and lesbians in terms of employment and housing, even to the point of serving as expert witnesses in court cases or helping in the writing of amicus curae briefs. They also argue against same sex marriage (Bradley, 2001; George, 2001).
With the rise of the gay liberation movement in the post-Stonewall era, overtly gay and lesbian perspectives began to be put forward in politics, philosophy and literary theory. Initially these often were overtly linked to feminist analyses of patriarchy (e.g., Rich, 1980) or other, earlier approaches to theory. Yet in the late 1980's and early 1990's queer theory was developed, although there are obviously important antecedents which make it difficult to date it precisely. There are a number of ways in which queer theory differed from earlier gay liberation theory, but an important initial difference can be gotten at by examining the reasons for opting for the term ‘queer’ as opposed to ‘gay and lesbian.’ Some versions of, for example, lesbian theory portrayed the essence of lesbian identity and sexuality in very specific terms: non-hierarchical, consensual, and, specifically in terms of sexuality, as not necessarily focused upon genitalia (e.g., Faderman, 1985). Lesbians arguing from this framework, for example, could very well criticize natural law theorists as inscribing into the very “law of nature” an essentially masculine sexuality, focused upon the genitals, penetration, and the status of the male orgasm (natural law theorists rarely mention female orgasms).
This approach, based upon characterizations of ‘lesbian’ and ‘gay’ identity and sexuality, however, suffered from three difficulties. First, it appeared even though the goal was to critique a heterosexist regime for its exclusion and marginalization of those whose sexuality is different, any specific or “essentialist” account of gay or lesbian sexuality had the same effect. Sticking with the example used above, of a specific conceptualization of lesbian identity, it denigrates women who are sexually and emotionally attracted to other women, yet who do not fit the description. Sado-masochists and butch/fem lesbians arguably do not fit this ideal of ‘equality’ offered. A second problem was that by placing such an emphasis upon the gender of one's sexual partner(s), other possible important sources of identity are marginalized, such as race and ethnicity. What is of utmost importance, for example, for a black lesbian is her lesbianism, rather than her race. Many gays and lesbians of color attacked this approach, accusing it of re-inscribing an essentially white identity into the heart of gay or lesbian identity (Jagose, 1996).
The third and final problem for the gay liberationist approach was that it often took this category of ‘identity’ itself as unproblematic and unhistorical. Such a view, however, largely because of arguments developed within poststructuralism, seemed increasingly untenable. The key figure in the attack upon identity as ahistorical is Michel Foucault. In a series of works he set out to analyze the history of sexuality from ancient Greece to the modern era (1980, 1985, 1986). Although the project was tragically cut short by his death in 1984, from complications arising from AIDS, Foucault articulated how profoundly understandings of sexuality can vary across time and space, and his arguments have proven very influential in gay and lesbian theorizing in general, and queer theory in particular (Spargo, 1999; Stychin, 2005).
One of the reasons for the historical review above is that it helps to give some background for understanding the claim that sexuality is socially constructed, rather than given by nature. Moreover, in order to not prejudge the issue of social constructionism versus essentialism, I avoided applying the term ‘homosexual’ to the ancient or medieval eras. In ancient Greece the gender of one's partner(s) was not important, but instead whether one took the active or passive role. In the medieval view, a ‘sodomite’ was a person who succumbed to temptation and engaged in certain non-procreative sex acts. Although the gender of the partner was more important than in the ancient view, the broader theological framework placed the emphasis upon a sin versus refraining-from-sin dichotomy. With the rise of the notion of ‘homosexuality’ in the modern era, a person is placed into a specific category even if one does not act upon those inclinations. What is the common, natural sexuality expressed across these three very different cultures? The social constructionist answer is that there is no ‘natural’ sexuality; all sexual understandings are constructed within and mediated by cultural understandings. The examples can be pushed much further by incorporating anthropological data outside of the Western tradition (Halperin, 1990; Greenberg, 1988). Yet even within the narrower context offered here, the differences between them are striking. The assumption in ancient Greece was that men (less is known about women) can respond erotically to either sex, and the vast majority of men who engaged in same-sex relationships were also married (or would later become married). Yet the contemporary understanding of homosexuality divides the sexual domain in two, heterosexual and homosexual, and most heterosexuals cannot respond erotically to their own sex.
In saying that sexuality is a social construct, these theorists are not saying that these understandings are not real. Since persons are also constructs of their culture (in this view), we are made into those categories. Hence today persons of course understand themselves as straight or gay (or perhaps bisexual), and it is very difficult to step outside of these categories, even once one comes to seem them as the historical constructs they are.
Gay and lesbian theory was thus faced with three significant problems, all of which involved difficulties with the notion of ‘identity.’ Queer theory thus arose in large part as an attempt to overcome them. How queer theory does so can be seen by looking at the term ‘queer’ itself. In contrast to gay or lesbian, ‘queer,’ it is argued, does not refer to an essence, whether of a sexual nature or not. Instead it is purely relational, standing as an undefined term that gets its meaning precisely by being that which is outside of the norm, however that norm itself may be defined. As one of the most articulate queer theorists puts it: “Queer is … whatever is at odds with the normal, the legitimate, the dominant. There is nothing in particular to which it necessarily refers. It is an identity without an essence” (Halperin, 1995, 62, original emphasis). By lacking any essence, queer does not marginalize those whose sexuality is outside of any gay or lesbian norm, such as sado-masochists. Since specific conceptualizations of sexuality are avoided, and hence not put at the center of any definition of queer, it allows more freedom for self-identification for, say, black lesbians to identify as much or more with their race (or any other trait, such as involvement in an S & M subculture) than with lesbianism. Finally, it incorporates the insights of poststructuralism about the difficulties in ascribing any essence or non-historical aspect to identity.
This central move by queer theorists, the claim that the categories through which identity is understood are all social constructs rather than given to us by nature, opens up a number of analytical possibilities. For example, queer theorists examine how fundamental notions of gender and sex which seem so natural and self-evident to persons in the modern West are in fact constructed and reinforced through everyday actions, and that this occurs in ways that privilege heterosexuality (Butler, 1990, 1993). Also examined are medical categories which are themselves socially constructed (Fausto-Sterling, 2000, is an erudite example of this, although she is not ultimately a queer theorist). Others examine how language and especially divisions between what is said and what is not said, corresponding to the dichotomy between ‘closeted’ and ‘out,’ especially in regards to the modern division of heterosexual/homosexual, structure much of modern thought. That is, it is argued that when we look at dichotomies such as natural/artificial, or masculine/feminine, we find in the background an implicit reliance upon a very recent, and arbitrary, understanding of the sexual world as split into two species (Sedgwick, 1990). The fluidity of categories created through queer theory even opens the possibility of new sorts of histories that examine previously silent types of affections and relationships (Carter, 2005).
Another critical perspective opened up by a queer approach, although certainly implicit in those just referred to, is especially important. Since most anti-gay and lesbian arguments rely upon the alleged naturalness of heterosexuality, queer theorists attempt to show how these categories are themselves deeply social constructs. An example helps to illustrate the approach. In an essay against gay marriage, chosen because it is very representative, James Q. Wilson (1996) contends that gay men have a “great tendency” to be promiscuous. In contrast, he puts forward loving, monogamous marriage as the natural condition of heterosexuality. Heterosexuality, in his argument, is an odd combination of something completely natural yet simultaneously endangered. One is born straight, yet this natural condition can be subverted by such things as the presence of gay couples, gay teachers, or even excessive talk about homosexuality. Wilson's argument requires a radical disjunction between heterosexuality and homosexuality. If gayness is radically different, it is legitimate to suppress it. Wilson has the courage to be forthright about this element of his argument; he comes out against “the political imposition of tolerance” towards gays and lesbians (Wilson, 1996, 35).
It is a common move in queer theory to bracket, at least temporarily, issues of truth and falsity (Halperin, 1995). Instead, the analysis focuses on the social function of discourse. Questions of who counts as an expert and why, and concerns about the effects of the expert's discourse are given equal status to questions of the verity of what is said. This approach reveals that hidden underneath Wilson's (and other anti-gay) work is an important epistemological move. Since heterosexuality is the natural condition, it is a place that is spoken from but not inquired into. In contrast, homosexuality is the aberration and hence it needs to be studied but it is not an authoritative place from which one can speak. By virtue of this heterosexual privilege, Wilson is allowed the voice of the impartial, fair-minded expert. Yet, as the history section above shows, there are striking discontinuities in understandings of sexuality, and this is true to the point that, according to queer theorists, we should not think of sexuality as having any particular nature at all. Through undoing our infatuation with any specific conception of sexuality, the queer theorist opens space for marginalized forms.
Queer theory, however, has been criticized in a myriad of ways (Jagose, 1996). One set of criticisms comes from theorists who are sympathetic to gay liberation conceived as a project of radical social change. An initial criticism is that precisely because ‘queer’ does not refer to any specific sexual status or gender object choice, for example Halperin (1995) allows that straight persons may be ‘queer,’ it robs gays and lesbians of the distinctiveness of what makes them marginal. It desexualizes identity, when the issue is precisely about a sexual identity (Jagose, 1996). A related criticism is that queer theory, since it refuses any essence or reference to standard ideas of normality, cannot make crucial distinctions. For example, queer theorists usually argue that one of the advantages of the term ‘queer’ is that it thereby includes transsexuals, sado-masochists, and other marginalized sexualities. How far does this extend? Is transgenerational sex (e.g., pedophilia) permissible? Are there any limits upon the forms of acceptable sado-masochism or fetishism? While some queer theorists specifically disallow pedophilia, it is an open question whether the theory has the resources to support such a distinction. Furthermore, some queer theorists overtly refuse to rule out pedophiles as ‘queer’ (Halperin, 1995, 62) Another criticism is that queer theory, in part because it typically has recourse to a very technical jargon, is written by a narrow elite for that narrow elite. It is therefore class biased and also, in practice, only really referred to at universities and colleges (Malinowitz, 1993).
Queer theory is also criticized by those who reject the desirability of radical social change. For example, centrist and conservative gays and lesbians have criticized a queer approach by arguing that it will be “disastrously counter-productive” (Bawer, 1996, xii). If ‘queer’ keeps its connotation of something perverse and at odds with mainstream society, which is precisely what most queer theorists want, it would seem to only validate the attacks upon gays and lesbians made by conservatives. Sullivan (1996) also criticizes queer theorists for relying upon Foucault's account of power, which he argues does not allow for meaningful resistance. It seems likely, however, that Sullivan's understanding of Foucault's notions of power and resistance are misguided.
The debates about homosexuality, in part because they often involve public policy and legal issues, tend to be sharply polarized. Those most concerned with homosexuality, positively or negatively, are also those most engaged, with natural law theorists arguing for gays and lesbians having a reduced legal status, and queer theorists engaged in critique and deconstruction of what they see as a heterosexist regime. Yet the two do not talk much to one another, but rather ignore or talk past one another. There are some theorists in the middle. For example, Michael Sandel takes an Aristotelian approach from which he argues that gay and lesbian relationships can realize the same goods that heterosexual relationships do (Sandel, 1995). He largely shares the account of important human goods that natural law theorists have, yet in his evaluation of the worth of same-sex relationships, he is clearly sympathetic to gay and lesbian concerns. Similarly, Bruce Bawer (1993) and Andrew Sullivan (1995) have written eloquent defenses of full legal equality for gays and lesbians, including marriage rights. Yet neither argue for any systematic reform of broader American culture or politics. In this they are essentially conservative. Therefore, rather unsurprisingly, these centrists are attacked from both sides. Sullivan, for example, has been criticized at length both by queer theorists (e.g., Phelan, 2001) and natural law theorists (e.g., George, 1999).
Yet as the foregoing also clearly shows, the policy and legal debates surrounding homosexuality involve fundamental issues of morality and justice. Perhaps most centrally of all, they cut to issues of personal identity and self-definition. Hence there is another, and even deeper, set of reasons for the polarization that marks these debates.
- Bawer, Bruce, 1993, A Place at the Table: The Gay Individual in American Society. New York: Poseidon Press.
- –––, 1996. Beyond Queer: Challenging Gay Left Orthodoxy. New York: The Free Press.
- Berman, Paul, 1993, “Democracy and Homosexuality” in The New Republic. Vol.209, No.25 (December 20): pp.17-35.
- Boswell, John, 1980, Christianity, Social Tolerance, and Homosexuality: Gay People in Western Europe from the Beginning of the Christian Era to the Fourteenth Century. Chicago: The University of Chicago Press.
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The United States has suspended food aid to North Korea in response to a planned rocket launch next month. Al Jazeera's Rosiland Jordan reports from Washington.
In international relations, aid (also known as international aid, overseas aid, or foreign aid) is a voluntary transfer of resources from one country to another, given at least partly with the objective of benefiting the recipient country. It may have other functions as well: it may be given as a signal of diplomatic approval, or to strengthen a military ally, to reward a government for behaviour desired by the donor, to extend the donor's cultural influence, to provide infrastructure needed by the donor for resource extraction from the recipient country, or to gain other kinds of commercial access. Humanitarianism and altruism are, nevertheless, significant motivations for the giving of aid.
Aid may be given by individuals, private organisations, or governments. Standards delimiting exactly the kinds of transfers that count as aid vary. For example, aid figures may or may not include transfers for military use: to cite one instance, the United States included military assistance in its aid figure until 1957 but no longer does. The most widely used measure of aid, "Official Development Assistance" (ODA) is such a figure. It is compiled by the Development Assistance Committee of the Organisation for Economic Co-operation and Development. The United Nations, the World Bank, and many scholars use the DAC's ODA figure as their main aid figure because it is easily available and reasonably consistently calculated over time and between countries. The DAC consists of 22 of the wealthiest Western industrialised countries plus the EU; it is a forum in which they coordinate their aid policies.
Humanitarian aid or emergency aid is rapid assistance given to people in immediate distress by individuals, organisations, or governments to relieve suffering, during and after man-made emergencies (like wars) and natural disasters. The term often carries an international connotation, but this is not always the case. It is often distinguished from development aid by being focussed on relieving suffering caused by natural disaster or conflict, rather than removing the root causes of poverty or vulnerability.
The provision of humanitarian aid or humanitarian response consists of the provision of vital services (such as food aid to prevent starvation) by aid agencies, and the provision of funding or in-kind services (like logistics or transport), usually through aid agencies or the government of the affected country. Humanitarian aid is distinguished from humanitarian intervention, which involves armed forces protecting civilians from violent oppression or genocide by state-supported actors.
The Geneva Conventions give a mandate to the International Committee of the Red Cross and other impartial humanitarian organizations to provide assistance and protection of civilians during times of war. The ICRC, has been given a special role by the Geneva Conventions with respect to the visiting and monitoring of prisoners of war.
The United Nations Office for the Coordination of Humanitarian Affairs (OCHA) is mandated to coordinate the international humanitarian response to a natural disaster or complex emergency acting on the basis of the United Nations General Assembly Resolution 46/182.
The Sphere Project handbook, Humanitarian Charter and Minimum Standards in Disaster Response, which was produced by a coalition of leading non-governmental humanitarian agencies, lists the following principles of humanitarian action:
Development aid is aid given by developed countries to support development in general which can be economic development or social development in developing countries. It is distinguished from humanitarian aid as being aimed at alleviating poverty in the long term, rather than alleviating suffering in the short term.
Official Development Assistance (ODA), mentioned above, is a commonly used measure of developmental aid. Development aid is given by governments through individual countries' international aid agencies and through multilateral institutions such as the World Bank, and by individuals through development charities such as ActionAid, Caritas, Care International or Oxfam.
The offer to give development aid has to be understood in the context of the Cold War. The speech in which Harry Truman announced the foundation of NATO is also a fundamental document of development policy:
In addition, we will provide military advice and equipment to free nations which will cooperate with us in the maintenance of peace and security. Fourth, we must embark on a bold new program for making the benefits of our scientific advances and industrial progress available for the improvement and growth of underdeveloped areas. More than half the people of the world are living in conditions approaching misery. Their food is inadequate. They are victims of disease. Their economic life is primitive and stagnant. Their poverty is a handicap and a threat both to them and to more prosperous areas. For the first time in history, humanity possesses the knowledge and skill to relieve the suffering of these people.
Aid from various sources can reach recipients through bilateral or multilateral delivery systems. Bilateral refers to government to government transfers. Multilateral institutions, such as the World Bank or UNICEF, pool aid from one or more sources and disperse it among many recipients.
Aid is often pledged at one point in time, but disbursements (financial transfers) might not arrive until later.
Private giving includes aid from charities, philanthropic organizations or businesses to recipient countries or programs within recipient counties.
Most monetary flows between nations are not counted as aid. These include market based flows such as foreign direct investments and portfolio investments. Remittances from migrant workers to their home countries are not counted in aid numbers even though, by volume, it is twice as large as what is counted as aid. Additionally, military support is not counted.
A major proportion of aid from donor nations is tied, mandating that a receiving nation spend on products and expertise originating only from the donor country. For example, Eritrea is forced to spend aid money on foreign goods and services to build a network of railways even though it is cheaper to use local expertise and resources. US law requires food aid be spent on buying food at home, instead of where the hungry live, and, as a result, half of what is spent is used on transport. Oxfam America and American Jewish World Service report that reforming US food aid programs could extend food aid to an additional 17.1 million people around the world.
The World Bank and the International Monetary Fund, as primary holders of developing countries' debt, attach structural adjustment conditionalities to loans which generally include the elimination of state subsidies and the privatization of state services. For example, the World Bank presses poor nations to eliminate subsidies for fertilizer even while many farmers cannot afford them at market prices. In the case of Malawi, almost five million of its 13 million people used to need emergency food aid. However, after the government changed policy and subsidies for fertilizer and seed were introduced, farmers produced record-breaking corn harvests in 2006 and 2007 as production leaped to 3.4 million in 2007 from 1.2 million in 2005, making Malawi a major food exporter. In the former Soviet states, the reconfiguration of public financing in their transition to a market economy called for reduced spending on health and education, sharply increasing poverty.
Aid is seldom given from motives of pure altruism; for instance it is often given as a means of supporting an ally in international politics. It may also be given with the intention of influencing the political process in the receiving nation. Whether one considers such aid helpful may depend on whether one agrees with the agenda being pursued by the donor nation in a particular case. During the conflict between communism and capitalism in the twentieth century, the champions of those ideologies, the Soviet Union and the United States, each used aid to influence the internal politics of other nations, and to support their weaker allies. Perhaps the most notable example was the Marshall Plan by which the United States, largely successfully, sought to pull European nations toward capitalism and away from communism. Aid to underdeveloped countries has sometimes been criticized as being more in the interest of the donor than the recipient, or even a form of neocolonialism.
S.K.B'. Asante lists some specific motives a donor may have for giving aid: defense support, market expansion, foreign investment, missionary enterprise, cultural extension. In recent decades, aid by organizations such as the International Monetary Fund and the World Bank has been criticized as being primarily a tool used to open new areas up to global capitalists, and being only secondarily, if at all, concerned with the wellbeing of the people in the recipient countries.
Besides criticism of motive, aid may be criticized simply on the grounds that it is not effective: i.e., it did not do what it was intended to do or help the people it was intended to help. This is essentially an economic criticism of aid. The two types of criticism are not entirely separate: critics of the ideology behind a piece of aid are likely to see it as ineffective; and indeed, ineffectiveness must imply some flaws in the ideology. Statistical studies have produced widely differing assessments of the correlation between aid and economic growth, and no firm consensus has emerged to suggest that foreign aid generally does boost growth. Some studies find a positive correlation, but others find either no correlation or a negative correlation. In the case of Africa, Asante (1985) gives the following assessment:
Summing up the experience of African countries both at the national and at the regional levels it is no exaggeration to suggest that, on balance, foreign assistance, especially foreign capitalism, has been somewhat deleterious to African development. It must be admitted, however, that the pattern of development is complex and the effect upon it of foreign assistance is still not clearly determined. But the limited evidence available suggests that the forms in which foreign resources have been extended to Africa over the past twenty-five years, insofar as they are concerned with economic development, are, to a great extent, counterproductive.
The economist William Easterly and others have argued that aid can often distort incentives in poor countries in various harmful ways. Aid can also involve inflows of money to poor countries that have some similarities to inflows of money from natural resources that provoke the resource curse.
Many individuals and organizations criticize U.S. Aid. Emergency funds from the International Monetary Fund (IMF) and World Bank, for instance, are linked to a wide range of free-market policy prescriptions that some argue interfere in a country's sovereignty. Policy prescriptions from outsiders can do more harm as they might not fit the local environment. The IMF can be good at helping countries over a short problematic financial period, but for poor countries with long lasting issues it can cause harm. In his book The White Man's Burden, Easterly argued that if the IMF only gave adjustment loans to countries that can repay it, instead of forgiving debts or lending repetitively even if conditions are not met, it would maintain its credibility.
In addition to the above criticisms, the logistics in which aid delivery occurs can be problematic. An earthquake in 2003 in Bam, Iran left tens of thousands of people in need of disaster zone aid. Although aid was flown in rapidly, regional belief systems, cultural backgrounds and even language seemed to have been omitted as a source of concern. Items such as religiously prohibited pork, and non-generic forms of medicine that lacked multilingual instructions came flooding in as relief. An implementation of aid can easily be problematic, causing more problems than it solves.
James Shikwati, a Kenyan economist, has argued that foreign aid causes harm to the recipient nations, specifically because aid is distributed by local politicians, finances the creation of corrupt government such as that led by Dr Fredrick Chiluba in Zambia bureaucracies, and hollows out the local economy. In an interview in Germany's Der Spiegel magazine, Shikwati uses the example of food aid delivered to Kenya in the form of a shipment of corn from America. Portions of the corn may be diverted by corrupt politicians to their own tribes, or sold on the black market at prices that undercut local food producers. Similarly, Kenyan recipients of donated Western clothing will not buy clothing from local tailors, putting the tailors out of business. In an episode of 20/20, John Stossel demonstrated the existence of secret government bank accounts which concealed foreign aid money destined for private purposes.
Some believe that aid is offset by other economic programs such as agricultural subsidies. Mark Malloch Brown, former head of the United Nations Development Program, estimated that farm subsidies cost poor countries about US$50 billion a year in lost agricultural exports:
"It is the extraordinary distortion of global trade, where the West spends $360 billion a year on protecting its agriculture with a network of subsidies and tariffs that costs developing countries about US$50 billion in potential lost agricultural exports. Fifty billion dollars is the equivalent of today's level of development assistance."
In response to aid critics, a movement to reform U.S. foreign aid has started to gain momentum. In the United States, leaders of this movement include the Center for Global Development, Oxfam America, the Brookings Institution, InterAction, and Bread for the World. The various organizations have united to call for a new Foreign Assistance Act, a national development strategy, and a new cabinet-level department for development.
As a result of these numerous criticisms, other proposals for supporting developing economies and poverty stricken societies. Some analysts, such as researchers at the Overseas Development Institute, argue that current support for the developing world suffers from a policy incoherence and that while some policies are designed to support the third world, other domestic policies undermine its impact, examples include:
One measure of this policy incoherence is the Commitment to Development Index (CDI) published by the Center for Global Development . The index measures and evaluates 22 of the world's richest countries on policies that affect developing countries, in addition to simply aid. It shows that development policy is more than just aid; it also takes into account trade, investment, migration, environment, security, and technology.
Thus, some states are beginning to go Beyond Aid and instead seek to ensure there is a policy coherence, for example see Common Agricultural Policy reform or Doha Development Round. This approach might see the nature of aid change from loans, debt cancellation, budget support etc., to supporting developing countries. This requires a strong political will, however, the results could potentially make aid far more effective and efficient.
It is true that aid is rarely given for motives of pure altruism. However, it is important to look at where aid goes. For example, “only about one fifth of U.S. aid goes to countries classified by the OECD as ‘least developed.’” This “pro-rich” trend is not unique to the United States. According to Collier, “the middle income countries get aid because they are of much more commercial and political interest than the tiny markets and powerlessness of the bottom billion.” What this means is that, at the most basic level, aid is not targeting the most extreme poverty.
The form of aid must also be considered. The World Bank, until recently, issued only loans, meaning that the country must repay both the loan and the interest rates. In contrast, the European Commission issues grants, which countries need not worry about paying back. This means that “loans have been going to the poorest countries and the grants to the middle-income countries.”
Furthermore, consider the breakdown, where aid goes and for what purposes. In 2002, total gross foreign aid to all developing countries was $76 billion. Dollars that do not contribute to a country’s ability to support basic needs interventions are subtracted. Subtract $6 billion for debt relief grants. Subtract $11 billion, which is the amount developing countries paid to developed nations in that year in the form of loan repayments. Next, subtract the aid given to middle income countries, $16 billion. The remainder, $43 billion, is the amount that developing countries received in 2002. But only $12 billion went to low-income countries ($15 billion for all developing countries) in a form that could be deemed budget support for basic needs.
The basic criticism of aid is that it neither goes where it was intended nor helps those intended. According to Collier, there are four known traps that contribute this problem. The first such trap is known as the conflict trap. Aid should not be used to finance military endeavors. It is difficult to “design aid in such a way that it works even in the environments of poor governance and poor policy that are most at risk of conflict.” The second trap is called the natural resource trap. Countries that are resource rich already have a large volume of capital flowing into their economies. However, it is not being used to its potential. The third trap occurs when a country is entirely landlocked, making it difficult for the country to engage in global trade.
The fourth trap is that of bad governance. However, “there are three ways in which aid can potentially help turnarounds: incentives, skills, and reinforcement.” Policy conditionalities, or structural adjustments, were reservations put on aid until a government agreed to aid implemented in the 1980s. This did not work. Aid needs to somehow provide incentives for giving the people power. Power needs to be transferred from the governments to the people. Aid should be restructured in order to allow for skills building in country. According to Collier, “technical assistance is not negligible – money spent on countries with the skilled people who constitute technical assistance is a quarter of total aid flows.” The problem is not that too little money is being provided, rather that technical assistance is not country specific. Aid is also given as budget support, reinforcement for failing states. There is an opportune moment for assisting failing states but it must be done at the right time. Aid cannot be continually poured into failing states and be expected to produce a turnaround. However if aid is given at the opportune political moment, it can support turnarounds. Collier suggests that when that moment occurs “pour in the technical assistance as quickly as possible to help implement reform” and “then, after a few years, start pouring in the money for the government to spend.”
Peter Singer argues that over the last three decades, “aid has added around one percentage point to the annual growth rate of the bottom billion.” He argues that this has made the difference between “stagnation and severe cumulative decline.” Aid can make progress towards reducing poverty worldwide, or at least help prevent cumulative decline.
Currently, donor institutions make proposals for aid packages to recipient countries. The recipient countries then make a plan for how to use the aid based on how much money has been given to them. Alternatively, NGO's receive funding from private sources or the government, and then implement plans to end their specific issues. In the views of many scholars, this system is inherently ineffective. If we hope to eliminate poverty, we must reexamine how we distribute funding, and how we attack problems.
According to Sachs, we should redefine how we think of aid. The first step should be to learn what developing countries hope to accomplish and how much money they need to accomplish those goals. Goals should be made with the Millennium Development Goals in mind for these furnish real metrics for providing basic needs. The “actual transfer of funds must be based on rigorous, country-specific plans that are developed through open and consultative processes, backed by good governance in the recipient countries, as well as careful planning and evaluation.”
Possibilities are also emerging as some developing countries are experiencing rapid economic growth, they are able to provide their own expertise gained from their recent transition. This knowledge transfer can be seen in donors, such as Brazil, whose $1 billion in aid outstrips that of many traditional donors. Brazil provides most of its aid in the form of technical expertise and knowledge transfers. This has been descried by some observers as a 'global model in waiting'.
Since the 1960s, improving the efficiency of foreign aid has been a common topic of academic research. There is debate on whether foreign aid is efficacious, but for the purposes of this article we will ignore that. Given that schema, a common debate is over which factors influence the overall economic efficiency of foreign aid. Indeed, there is debate about whether aid impact should be measured empirically at all, but again, we will limit our scope to increasing the economic efficiency.
At the forefront of the aid debate has been the conflict between professor William Easterly of New York University and his ideological opposite, Jeffrey Sachs, from Columbia University. Easterly advocates the "searcher's" approach, while Sachs advocates a more top down, broad planned approach. We will discuss both of these at length.
William Easterly offers a nontraditional, and somewhat controversial “searching” approach to solving poverty, as opposed to the “planned” approach in his famous critique of the more traditional Owen/Sachs, The White Man’s Burden. Traditional poverty reduction, Easterly claims is based on the idea that we know what is best for those countries, which are impoverished. He claims that they know what’s best. Having a top down “master plan,” he claims, is inefficient. His alternative, called the “Searchers” approach, uses a bottom up strategy. That is, this approach starts by surveying the poor in the countries in question, and then tries to directly aid individuals, rather than governments. Local markets are a key incentive structure. The primary example is of Mosquito nets in Malawi. In this example, an NGO sells Mosquito nets to rich Malawians, and uses the profits to subsidize cheap sales to the impoverished. Hospital nurses are used as middle-women, profiting a few cents on every net sold to a patient. This incentive structure has seen the usage of nets in Malawi spike over 40% in less than 7 years.
One of the central tenets in Easterly’s approach is a more bottom up philosophy of aid. This applies not only to the identification of problems, but to the actual distribution of capital to the areas in need. In effect, Easterly would have countries go to the area which needed aid, collect information about the problem, find out what the population wanted, and then work from there. In keeping with this, funds would also be distributed from the bottom up, rather than being given to a specific government.
Easterly also advocates working through currently existing Aid organizations, and letting them compete for funding. Utilizing preexisting national organizations and local frameworks would not only help give target populations a voice in implementation and goal setting, but is more efficient economically. Easterly argues that the preexisting frameworks already "know" what the problems are, as opposed to outside NGOs who tend to "guess".
Easterly strongly discourages aid to government as a rule. He believes, for several reasons, that aid to small “bottom up” organizations and Individual groups is a better philosophy than to large governments.
Easterly states that for far too long, inefficient aid organizations have been funded, and that this is a problem. The current system of evaluation for most aid organizations is internal. Easterly claims that the process is biased because organizations have a large incentive to represent their progress in a positive light. What he proposes as an alternative is an independent auditing system for aid organizations. Before receiving funding, the organization would state their goals and how they expect to measure and achieve them. If they do not meet their goals, Easterly proposes we shift our funding to organizations who are successful. This would prompt organizations to either become efficient, or obsolete.
Easterly believes that aid goals should be small. In his opinion, one of the main failings of aid lies in the fact that we create large, utopian lists of things we hope to accomplish, without the means to actually see them to fruition. Rather than establish a utopian vision for a particular country, Easterly insists that we shift our focus to the most basic needs and improvements.If we feed, clothe, vaccinate, build infrastructure, and support markets, the macroscopic results will follow.
The “Searching Approach” is intrinsically tied to the market. Easterly claims that the only way for poverty to truly end is for the poor to be given the capability to lift themselves out of poverty, and then for it to happen. Philosophically, this sounds like the traditional “bootstrap” theory, but it isn’t. What he says is that the poor should be given the fiscal support to create their market, which would give them the ability to become self reliant in the future.
In the end of his book, Easterly proposes a voucher system for foreign aid. The poor would be distributed a certain amount of vouchers, which would act as currency, redeemable to aid organizations for services, medicines, and the like. These vouchers would then be redeemed by the aid organizations for more funding. In this way, the aid organization would be forced to compete, if by proxy.
Sachs presents a near dichotomy to Easterly. Sachs presents a broad, proscriptive solution to poverty. In his book, The End of Poverty, he explains how throughout history, countries have ascended from poverty by following a relatively simple model. First, you promote agricultural development, then industrialize, embrace technology, and finally become modern. This is the standard “western” model of development that has been followed by countries such as China and India. Sachs main idea is that there should have a broad analytical “checklist” of things a country must attain before it can reach the next step on the ladder to development. Western nations should donate a percentage of their GDP as determined by the UN, and pump money into helping impoverished countries climb the ladder. Sachs insists that if followed, his strategy would eliminate poverty by 2025.
Sachs advocates using a top down methodology, utilizing broad ranging plans developed by external aid organizations like the UN and World Bank. To Sachs, these plans are essential to a coherent and timely eradication of poverty.He surmises that if donor and recipient countries follow the plan, they will be able to climb out of poverty.
Part of Sachs’ philosophy includes strengthening the International Monetary Fund, World Bank, and the United Nations. If those institutions are given the power to enact change, and freed from mitigating influences, then they will be much more effective. Sachs does not find fault in the international organizations themselves. Instead, he blames the member nations who compose them. The powerful nations of the world must make a commitment to end poverty, then stick to it.
Sachs believes it is best to empower countries by utilizing their existing governments, rather than trying to circumnavigate them. He remarks that while the corruption argument is logically valid in that corruption harms the efficiency of aid, levels of corruption tend to be much higher on average for countries with low levels of GDP. He contends that this hurdle in government should not disqualify entire populations for much needed aid from the west.
Sachs does not see the need for independent evaluators, and sees them as a detractor to proper progress. He argues that many facets of aid cannot be effectively quantified, and thus it is not fair to try to put empirical benchmarks on the effectiveness of aid.
Sachs’ view makes it a point to attack and attempt to disprove many of the ideas that the more “pessimistic” Easterly stands on.
First, he points to economic freedom. One of the common threads of logic In aid is that countries need to develop economically in order to rise from Poverty. On this, there is not a ton of debate. However, Sachs contends that Easterly, and many other neo-Liberal economists believe high levels of Economic freedom in these emerging markets is almost a necessity to Development. Sachs himself does not believe this. He cites the lack of Correlation between the average degrees of Economic Freedom in countries And their yearly GDP growth, which in his data set is completely Inconclusive.
Also, Sachs contends that democratization is not an integral part of having Efficient aid distribution. Rather than attach strings to our aid dollars, or only Working with democracies or “good governments”, Sachs believes we should Consider the type of government in the needy country as a secondary Concern.
Sachs’ entire approach stands on the assertion that abject poverty could be Ended worldwide by 2025.
Dollar/Collier showed that current allocations of aid are allocated inefficiently. They came to the conclusion that aid money is given in many cases as an incentive to change policy, and for political reasons, which in many cases can be less efficient than the optimal condition. They agree that bad policy is detrimental to economic growth, which is a key component of poverty reduction, but have found that aid dollars do not significantly incentivize governments to change policy. In fact, they have negligible impact. As an alternative, Dollar proposes that aid be funneled more towards countries with “good” policy and less than optimal amounts of aid for their massive amounts of poverty. With respect to “optimal amounts” Dollar calculated the marginal productivity of each additional dollar of foreign aid for the countries sampled, and saw that some countries had very high rates of marginal productivity (each dollar went further), while others [with particularly high amounts of aid, and lower levels of poverty] had low [and sometimes negative] levels of marginal productivity. In terms of economic efficiency, aid funding would be best allocated towards countries whose marginal productivities per dollar were highest, and away from those countries who had low to negative marginal productivities. The conclusion was that while an estimated 10 million people are lifted from poverty with current aid policies, that number could be increased to 19 million with efficient aid allocation.
In a followup study to Collier, Dalgaard finds that regardless of policies and utility, aid seems to be less efficient per dollar in the tropics. He states that aid dollars are likely to be just as effective anywhere, so there must be a mitigating externality that is not being accounted for.
New Conditionality is the term used in a paper to describe somewhat of a compromise between Dollar and Hansen. Paul Mosely describes how policy is important, and that aid distribution is improper. However, unlike Dollar, “New Conditionality” claims that the most important factors in efficiency of aid are income distributions in the recipient country and corruption.
One of the problems in foreign aid allocation is the marginalization of the fragile state. The fragile state, with its high volatility, and risk of failure scares away donors. The people of those states feel harm and are marginalized as a result. Additionally, the fate of neighboring states is important, as economies of the directly adjacent states to those impoverished, volatile “fragile states” can be negatively impacted by as much as 1.6% of their GDP per year. This is no small figure. McGillivray advocates for the reduced volatility of aid flows, which can only be attained through analysis and coordination.
Beynon revisits Collier/Dollar’s “efficient aid allocation” and finds much of the same results, although he believes that the model needs to be more robust, with more samplings, and take into account other Millennium Development Goals.
Aid often comes with conditions on its allocation. This is called Aid Conditionality. The two types of conditionality are Outcome and Process based conditionality. Outcome conditionality ties aid to a certain goal, while process conditionality ties aid to a certain method of implementation. These two forms of conditionality are often seen together. The problem with aid conditionality is that it not only restricts the local legislatures in how they shape their own country, but often removes local populations from the goal setting and decision making process entirely.
Academic research has shown that in many instances, aid is conditionally tied due to political motives, rather than notions of proper policy and implementation.
The European Union double standard of trade policy mentioned above is more illustrated in its trade with Latin America. In their April 2002 publication,Oxfarm Report reveals that aid tied to trade Liberalization by the donor countries such as the EU with the aim of achieving economic objective is becoming detrimental to developing countries.For example, the EU subsidizes its agricultural sectors in the expense of Latin America who must liberalize trade in order to qualify for aid. Latin America having, a comparative advantage on agriculture and mainly forms its export sector losts $ 4bn annually due to the EU farm subsidy policies. The aid money actually, instead of being used to finance other development infrastructure, ends up being used to fund terms of trade deficit brought about by trade liberalization. This inequalities and many others imposed upon developing countries have actually benefited the donor rather than the recipient.This cases clearly illustrate the evolution of development aid assistance from providing financial resource to fuel development to being a foreign policy tool used by donor countries to promote their domestic interests. while this can be considered controversial issues, addressing them adequately may help attain some effectiveness in development aid assistance goals.
Axel Dreher showed that over the period 1951-2004, there was a robust positive relationship between United Nations Security Council membership and participation in Internation Monetary Fund programs. Additionally, conditions placed on Security Council member nations' aid are lower in number and severity. Member nations of the UNSC received on average .8 more World Bank loans per year.
Seemingly in support of Bill Easterly's searcher's approach, Carlos Santiso found evidence that conditionality is quite harmful to development, and advocates a "radical approach in which donors cede control to the recipient country".
Santiso also argues that neither good governance or democracy is sustainable without the other.
Former USAID official Carol Lancaster, in her book Foreign Aid (2007) defines foreign aid as: "a voluntary transfer of public resources, from a government to another independent government, to an NGO, or to an international organization (such as the World Bank or the UN Development Program) with at least a 25 percent grant element, one goal of which is to better the human condition in the country receiving the aid." (p 9.)
Both definitions employ the concept that benefit to the people of the receiving country must be one but not necessarily the only objective.
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Category Archives: War Crimes and Criminals
Current British premier David Cameron praised Lady Thatcher for having “saved Britain” and for making the has-been colonial power “great again”.
Tributes poured forth from French and German leaders, Francoise Hollande and Angela Merkel, while US President Barack Obama said America had lost a “special friend”.
Former American secretary of state Henry Kissinger and former Russian leader Mikhail Gorbachev also lamented the loss of “an historic world figure”. Polish ex-president Lech Walesa hailed Margaret Thatcher for having brought down the Soviet Union and Communism.
Such fulsome praise may be expected coming from so many war criminals. But it is instructive of how history is written by the victors and criminals in high office. Obama, Cameron, Hollande and Merkel should all be arraigned and prosecuted for war crimes in Iran, Iraq, Afghanistan, Libya, Syria, Pakistan, Somalia and Mali, among other places. Kissinger has long evaded justice for over four decades for his role in the US genocide in Southeast Asia during the so-called Vietnam War in which over three million people were obliterated in Vietnam, Laos and Cambodia.
The British state is to give Thatcher, who died this week aged 87, a full military-honours funeral. The praise, eulogies, wreaths and ceremonies are all self-indictments of association with one of the most ruthless and criminal political figures in modern times.
So, here is a people’s history of Thatcher’s legacy.
She will be remembered for colluding with the most reactionary elements of Rupert Murdoch’s squalid media empire to launch a war over the Malvinas Islands in 1982, a war that caused hundreds of lives and involved the gratuitous sinking of an Argentine warship, the Belgrano,
by a British submarine.
By declaring war, rather than conducting political negotiations with Argentina over Britain’s ongoing colonial possession of the Malvinas, Thatcher salvaged her waning public support in Britain, and the bloodletting helped catapult her into a second term of office in Downing Street. Her political “greatness” that so many Western leaders now eulogize was therefore paid in part by the lives of Argentine and British soldiers, and by bequeathing an ongoing source of conflict in the South Atlantic.
It wasn’t just foreigners that Thatcher declared war on. Armed with her snake-oil economic policies of privatisation, deregulation, unleashing finance capitalism, pump-priming the rich with tax awards subsidised by the ordinary working population, Thatcher declared war on the British people themselves. She famously proclaimed that “there was no such thing as society” and went on to oversee an explosion in the gap between rich and poor and the demolition of social conditions in Britain. That legacy has been amplified by both successive Conservative and Labour governments and is central to today’s social meltdown in Britain – more than two decades after Thatcher resigned. Laughably, David Cameron, a protégé of Thatcher, claims that she “saved” Britain. The truth is Thatcher accelerated the sinking of British capitalism and society at large. What she ordered for the Belgrano has in a very real way come to be realised for British society at large.
During her second term of office in the mid-1980s, the Iron Lady declared war on the “enemy within”. She was referring to Britain’s strongly unionised coal-mining industry. Imagine declaring war on your own population. That is a measure of her pathological intolerance towards others who did not happen to share her obnoxious ideological views – ideological views that have since become exposed as intellectually and morally bankrupt.
For over a year around 1984, her Orwellian mindset and policies starved mining communities in the North of England into submission. Her use of paramilitary police violence also broke the resolve and legitimate rights of these communities. Miners’ leader Arthur Scargill would later be vindicated in the eyes of ordinary people, if not in the eyes of the mainstream media. Britain’s coalmines were systematically shut down, thousands of workers would be made unemployed, and entire communities were thrown on the social scrap heap. All this violence and misery was the price for Thatcher’s ideological war against working people and their political rights.
The class war that Thatcher unleashed in Britain is still raging. The rich have become richer, the poor decidedly more numerous and poorer. The decimation of workers’ rights and the unfettered power given to finance capital were hallmarks of Thatcher’s legacy and are to this day hallmarks of Britain’s current social decay. But that destructive legacy goes well beyond Britain. The rightwing nihilistic capitalism that Thatcher gave vent to was and became a zeitgeist for North America, Europe and globally. The economic malaise that is currently plaguing the world can be traced directly to such ideologues as Margaret Thatcher and former US President Ronald Reagan.
A final word on Thatcher’s real legacy, as opposed to the fakery from fellow war criminals, is her role in Ireland’s conflict. Her epitaph of “Iron Lady” is often said with admiration or even sneaking regard for her supposed virtues of determination and strength. In truth, her “iron” character was simply malevolent, as can be seen from her policies towards the Irish struggle for independence from Britain. In 1981, 10 Irish republican prisoners, led by a young Belfast man by the name of Bobby Sands, died from hunger strikes. The men died after more than 50 days of refusing prison food because they were demanding to be treated as political prisoners, not as criminals. Thatcher refused to yield to their demands, denouncing them as criminals and callously claiming that they “took their own lives”. No matter that Bobby Sands had been elected by tens of thousands of Irish voters to the British House of Parliament during his hunger strike. He was merely a criminal who deserved to die, according to the cold, unfeeling Thatcher.
As a result of Thatcher’s intransigence to negotiate Irish rights, the violence in the North of Ireland would escalate over the next decade, claiming thousands of lives. As with Las Malvinas dispute with Argentina, Thatcher deliberately took the military option and, with that, countless lives, rather than engage in reasoned, mutual dialogue. Her arrogance and obduracy blinded her to any other possibility.
As the violence gyrated in Ireland, Thatcher would also embrace the criminal policy of colluding with pro-British death squads. Armed,funded and directed by British intelligence, these death squads would in subsequent years kill hundreds of innocent people – with the knowledge and tacit approval of Lady Thatcher. It was a policy of British state terrorism in action, sanctioned by Thatcher. One of those victims was Belfast lawyer Pat Finucane, who was murdered in February 1989. He was shot 12 times in the head in front of his wife and children by a British death squad, after the killers smashed their way into the Finucane home on a Sunday afternoon.
Thus whether in her dealings with Las Malvinas row with Argentina, the British working people or Irish republicans, Margaret Thatcher was an intolerant militarist who always resorted to demagoguery, violence and starvation to get her political way. She was a criminal fascist who is
now proclaimed to be a national hero.
Reports this week say that Thatcher died with Alzheimer’s, the brain-degenerating disease in which the patient loses their faculty for memory. Western leaders, it seems, would also like to erase public memory of Thatcher’s criminal legacy.
Jewish hasbara bodies aiding and abetting Israeli crimes must face the full force of the law
By Nureddin Sabir
Editor, Redress Information & Analysis
The world of Jewish politics is so back to front and upside down that, when it comes to Israel, bad is good and wrong is right.
That is the sad fact of which we need to remind Israel flag wavers, such as failed US politician Katrina Lantos Swett, who from time to time rear their heads to bleat “anti-Semitism” and decry the “deligitimization of Israel”, which they blame for allegedly rising anti-Jewish sentiment.
On 5 April the Times of Israel reported that a Jewish organization in Melbourne, Australia, could face expulsion from the country’s Jewish umbrella body for launching a campaign that calls for the boycott of products from West Bank settlements.
These settlements – colonies and squatter camps, in fact – are illegal under international law and, with their rapid expansion under successive Israeli governments, both of the left and the right, are terminally undermining any prospect of a two-state solution of the Palestinian-Israeli conflict.
Boycott settlements call
The Australian campaign “Don’t Buy from the Settlements” was launched on 26 March by the Melbourne-based Australian Jewish Democratic Society (AJDS) with the aim of encouraging Australian Jews to avoid buying products made in Jewish colonies and squatter camps located in the Palestinian territories occupied after the 1967 war.
In a media release published on its website, the AJDS said:
Israeli settlements are seen around the world as a major obstacle to creating peace between Israelis and Palestinians. One way to take a stand against the harm they create is not to buy the products they produce. This sends a clear message that we will not be complicit in the settlement programme.
According to Jordy Silverstein, an executive member of the AJDS:
Not buying products from settlements will not work on its own, but it is one small step that we can take. When we add in the possibility of sharing knowledge about what the settlements mean and what they do … we can work alongside Palestinians, Israelis and people throughout the diasporas to create an exciting, liberating future.
However, this principled stance of the AJDS quickly prompted Australia’s Israel hasbara (propaganda) groups to gang up against it. For them, anything short of total, unconditional support for Israel, its illegal occupation and colonization of Arab territories, its war crimes and its crimes against humanity is tantamount to treason.
Orgy of Jewish solidarity with Israel
So, first came the umbrella Jewish organization, the Zionist Federation of Australia, which described the call for boycott as “immoral” and “repugnant”.
Then came the Jewish Community Council of Victoria, which consists of more than 50 Jewish organizations in the state. Its president, Nina Bassat, touted the idea of expelling the AJDS for having the audacity to call for action against Israeli criminality.
Bassat was joined in this orgy of Jewish tribal solidarity with the racist Israeli state by Peter Wertheim, of the Executive Council of Australian Jewry, who also called for the AJDS to be expelled from Australia’s Jewish umbrella body.
“The AJDS campaign is repugnant to the strong anti-BDS policies of every Jewish communal roof body in Australia,” he bleated, referring to the boycott, divestment and sanctions movement, “and to the ECAJ [Executive Council of Australian Jewry] platform of support for Israel and its legitimacy as the state of the Jewish people”.
“Anti-Semitism” and support for Israel
If there is indeed a rise in anti-Jewish sentiment anywhere in the world, then to find the cause look no further than the above-mentioned types of frenzied Jewish defence of Israeli crimes.
There is a strong case for nurturing a worldwide culture at grassroots level that promotes international legality and opposes crimes committed or sponsored by states. As a first step towards this, it is necessary to bring to account groups and individuals that support criminal actions committed or sponsored by states.
Just as in most civilized countries there are laws against aiding, abetting and glorifying terrorism, and against crimes such as rape and drugs trafficking, so there should also be laws that proscribe supporting or glorifying state crimes.
If neo-Nazis and Nazi holocaust deniers can be prosecuted in Germany, France and elsewhere in Europe, so should organizations that support Israeli crimes also be prosecuted, whether in Australia or anywhere else, for aiding, abetting and/or glorifying Israeli war crimes and crimes against humanity.
During the presentation by Ivanov data was shared which shows that 820 tons of heroin makes its way into Europe and Russia every year. Most of it is trafficked through the unstable Middle East and Africa as is part of the 994 tons of cocaine that is consumed in the U.S. and Europe.
“… First, there is no impartial intervention. Entering the conflict is to take sides. Ronald Reagan, 241 Marines, and 17 American embassy personnel learned that lesson in Lebanon in 1983. Washington had proclaimed its commitment to peace by aiding one force in a multi-sided civil war. By becoming a de facto combatant the administration turned Americans into targets. Aiding Syria’s opposition means becoming a participant in that conflict.
Paradoxically, aiding the resistance could drive some Syrians who desire a negotiated solution toward the government. The Financial Times recently reported: “As the civil war becomes ever dirtier, rebels’ actions are starting to mirror those of the regime.” In fact, opposition fighters increasingly kill regime soldiers and supporters, and have turned to crime, including kidnapping, to raise funds.
Second, there is no magic elixir that combines riskless intervention with speedy conquest. In Libya the allies provided the rebels with air support, but only enough to drag out the conflict for five months, during which time thousands, and perhaps tens of thousands, of Libyans died. By being prudent and cost-conscious the allies were not humanitarian, their professed objective….”
The 9/11 attacks were criminal acts intended as reprisals for American-sponsored oppression in the Middle East
Anti-war campaigners gathered in Trafalgar Sqaure on April 1, 2009 to protest against the ongoing wars in Afghanistan and Iraq. The protest, organised by the Stop the War Coalition and the Campaign for Nuclear Disarmament, was part of a day of protests in advance of tomorrows G20 summit in London. The protest began outside the US Embassy and then marched through Central London to Trafalgar Square. (Photo: Simon Kimber)
River to Sea Uprooted Palestinian
The views expressed in this article are the sole responsibility of the author and do not necessarily reflect those of this Blog!
The U.S. Global Strangulation of Iran has nothing to do in the short or intermediate terms with any fear Iran will get a nuclear weapon. American intelligence and former heads of Israel’s Mossad agree with the international assessment that there is no evidence Iran is now trying to produce nuclear weapons or that they could do so in less than several to ten years if they tried to do so. There is also nearly universal agreement among serious analysts, including U.S. intelligence sources going back many years, that attacking Iran would likely lead to severe cuts in oil supplies from the Persian Gulf that would produce a severe shock to the global economy, especially in this time of financial crisis. This does not mean Obama or the entire U.S. leadership at the top is going insane or trying to destroy the U.S. Empire. There is a “method to this seeming madness,” a very clear grand strategic goal involved in these constant threats and attacks on Iran. The general goal is to save the U.S. Empire and its domination and exploitation of the world to stop the accelerating decline and disintegration of the Empire and the American financial system and economy. The specific goal of all the attacks on Iran is to stop the rise of the Iranian-Shia superpower in the Middle East and its rapidly growing alliance with Islamist Sunni powers growing across North Africa, the Middle East, South Central and Central Asia. Iran is also the core leader in the growing alliance of Shia and Islamist Sunnis with the vast nations of Russia, India and China against the U.S. Empire’s energy and dollar imperialism. all of this makes Iran and its allies a mortal threat to the U.S. global empire and the present U.S. financial system tied to that empire. this is why the U.S. is taking such desperate risks to try to strangle Iran and to destroy its central position in this anti-American imperialism coalition that is growing rapidly. Iran and its Shia Allies across the whole Middle East are a very real threat to the U.S. Global Empire’s domination and exploitation of this whole vast region. That Iranian superpower Shia bloc’s rapidly growing alliance with Islamist Sunnis makes Iran the core of a very grave threat to the U.S. global empire which I believe has already developed far enough to prevent the U.S. from continuing its imperialist suppression of muslim nations with totalitarian, puppet regimes and the use of those to exploit this vast region. But I think these would be temporary set-backs for Iran and its great alliance of nations lining up against the U.S. Empire. I think the world sees the Giant is crippled severely by its losing wars, by the decline of the U.S. economy relative to the rising nations around the world, the so-called BRICKS, though one must add another I to that for Iran [in addition to India] – BRICKS, and The Great Global Financial Crisis the U.S. produced and is making worse and worse with vast help from its puppets in the EU and elsewhere, as even the “liberal” Brookings Institute just reported on the basis of global indices they analyze.
By attacking Iran and all of the Shia and Islamist Sunni nations allying more and more closely with Iran, most obviously Syria right now, the U.S. is in fact turning the vast Muslim World of nearly one and a half billion mostly young and very fertile people and much of the rest of the world against its brutal Empire.
This desperate grand strategy has, therefore, put the existence of the U.S. empire and of the U.S. financial system dependent on it at grave risk. at the very least, the U.S. is “going for broke” – “playing for all the marbles” – in its blatantly brutal, imperialistic strangulation of Iran and now some of its allies such as Syria.
The U.S. destroyed the Sunni block to the rise of Shia Iran as a Super Power in the Middle East when it destroyed Hussein’s Iraqi military machine, which the U.S. had earlier secretly supported in Iraq’s vast invasion of Iran [which led to the defeat of Iraq and the U.S.] Sunni Iraq was the only major block to the rise of the Shia as a Super Power at the heart of the Middle East. When the U.S. also installed a totalitarian Shia military power in Iraq which from the beginning was quietly allied with their brother Shia in Iran [where many of the new rulers of Iraq had lived for many years to escape Hussein's secret police], the new, rising Shia Super Alliance of Iran and Iraq became a huge military power of 100,000,000 mostly young people who loathe the U.S. and control more oil and gas than the Saudi Arabian and UAE puppets of the U.S. Iran was by then already building its Super Alliance with Syria, Hizbollah in Lebanon and Hamas in Palestine, an Alliance at the very heart of the Arab World and the Middle East. That alliance grew very rapidly after Hizbollah defeated Israel and drove them out of Lebanon for the second time.
Beginning over a year ago the U.S. puppet Sunni regimes across North Africa started falling to popular revolutions against the puppets and, above all, against the U.S. Empire which replaced the U.K. and French Empires that had oppressed them earlier. The revolutions are still going on in almost every U.S. Sunni puppet nation, but in almost all of them the people are winning their revolutions and moving steadily into a new age of Islamist governments which are in fact very anti-American and becoming more so as the U.S. allys itself more and more openly with the most hated totalitarian regimes, such as the Egyptian Army. the absurd, monarchical and pseudo-democratic tyrannies in Saudi Arabia, the UAE, Yemen and a few others.
From the beginning of these new Sunni Islamist movements and governments, the Iranians and their Shia allies have allied themselves more and more with the Sunni Islamists, based largely on their common loathing of the Great Satan – the U.S. – and the terrorist threat it poses to all of them. Turkey has not had an Islamist revolution, but the Sunni Islamists are taking over more and more in Turkey and becoming more and more Islamist, more and more anti-American and more and more quietly allied with the Iranians and the rising Sunni Islamists.
At the same time the U.S. was enraging the people of the Middle East against itself, the U.S. was enraging the Afghans and Pakistanis in the same way only more virulently with its terrorist murders and raids everywhere. The hearts and minds of nearly 200,000,000 Afghans and Pakistanis are now virulently anti-American. Though they are not yet allied with the rising Iranian Super Alliance, they are becoming more friend;y and will almost certainly find more and more common causes against the U.S. as the U.S. and its European [Nato] puppets flee the rising power of the Islamist guerilla armies. The U.S. has lost its terrorist war against Afghanistan and Pakistan and pushed them closer and closer to both iran and china, and even to some degree india and russia, to counter-balance the u.s. in its final desperate war against the world.
All of these vast developments are being hidden from the American and European peoples as much as possible by vast propaganda Media Wars against Iran and its grand global alliance against the U.S. But intelligent leaders around the world see the U.S. is desperate and is losing and even in Latin America the few remaining puppet regimes of the U.S. Empire, notably Mexico and Colombia, are turning more and more against the U .S. even in full public view. This happened just this weekend in Cartagena, Colombia, where Obama and his motley crowd of Key Stone Cop “security guards” were publicly humiliated by the still partially puppet regime of Colombia which needs vast U.S. military help in oppressing the peasants of the North for the vastly rich landowning families of the South. All 30 Latin Leaders there opposed even a future meeting of this sort with the U.S. without the U.S. allowing Cuba to attend and almost all insist as well on an end for support of the U.K. puppet regime’s claims to the Falklands off the coast of Argentina. Hugo Chavez, the elected president of Venezuela, was not even there because he is getting more cancer treatment in Cuba [not Miami]. He is more and more an ally of Iran, as is even Brazil on crucial matters like opposition to Dollar Imperialism and Inflation Thievery. Had he and his few very close allies from Latin America, Ecuador and Nicaragua, been there, the anti-American thrust would have been even greater. As it is, the Cartagena police wound up mortally embarrassing the Secret Key Stone Cops, the U.S. Army and Obama by taking U.S. security guards into custody over refusal to pay a prostitution bill for servicing rendered. Obama et al were made a laughing stock in a way that never happened before when these terrorized nations were ruled by oppressed puppet regimes. [Like Congressman Issa, I think these were SOP activities which led to police custody and massive public exposure for the first time. The U.S. Empire is losing its grip even on still dependent puppets.]
The U.S. has gotten very desperate in the Middle East. It is not only risking a cut off of oil by Iran over the U.S. strangulation of Iran, but the U.S. is now more and more openly involved in supporting the Sunni revolt against Syria and the Sunni oppression of Shia revolt against Bahrain by Saudia Arabia, the only major U.S. puppet left in the Middle East and a tyrannical monarchy loathed by the whole vast region and its own population [as far as we can tell in such an extremely closed police state].
It is always possible the U.S. will realize it is risking its entire Empire in this grand strategy, or that Iran will buckle under the strangulation pressures and retreat, or an Iranian ally like Syria will start a war inadvertently with another, more important ally, Sunni Turkey, that might reverse the growing great alliance. Maybe the U.S. will become more wise in its mix of sticks and carrots in dealing with Iran.
I think the U.S. is desperately fighting its last great imperial war against the world. I think the U.S. will lose this WAR and destroy its Empire and its global financial system. Obama has become the Undertaker of the U.S. Global Empire and, if he and his desperate strategists keep attacking the world, he will be the Undertaker of the U.S. as it now exists. By all indications, they will continue to destroy the U.S. unless some greater power, perhaps a Great Awakening of the American people who would prefer not to be impoverished or torn apart by the Obama Undertaker Strategy.
River to Sea Uprooted Palestinian
The views expressed in this article are the sole responsibility of the author and do not necessarily reflect those of this Blog!
The U.S. Global Strangulation of Iran has nothing to do in the short or intermediate terms with any fear Iran will get a nuclear weapon. American intelligence and former heads of Israel’s Mossad agree with the international assessment that there is no evidence Iran is now trying to produce nuclear weapons or that they could do so in less than several to ten years if they tried to do so. There is also nearly universal agreement among serious analysts, including U.S. intelligence sources going back many years, that attacking Iran would likely lead to severe cuts in oil supplies from the Persian Gulf that would produce a severe shock to the global economy, especially in this time of financial crisis.
This does not mean Obama or the entire U.S. leadership at the top is going insane or trying to destroy the U.S. Empire. There is a “method to this seeming madness,” a very clear grand strategic goal involved in these constant threats and attacks on Iran. The general goal is to save the U.S. Empire and its domination and exploitation of the world to stop the accelerating decline and disintegration of the Empire and the American financial system and economy. The specific goal of all the attacks on Iran is to stop the rise of the Iranian-Shia superpower in the Middle East and its rapidly growing alliance with Islamist Sunni powers growing across North Africa, the Middle East, South Central and Central Asia.
Iran is also the core leader in the growing alliance of Shia and Islamist Sunnis with the vast nations of Russia, India and China against the U.S. Empire’s energy and dollar imperialism.
all of this makes Iran and its allies a mortal threat to the U.S. global empire and the present U.S. financial system tied to that empire. this is why the U.S. is taking such desperate risks to try to strangle Iran and to destroy its central position in this anti-American imperialism coalition that is growing rapidly.
Iran and its Shia Allies across the whole Middle East are a very real threat to the U.S. Global Empire’s domination and exploitation of this whole vast region. That Iranian superpower Shia bloc’s rapidly growing alliance with Islamist Sunnis makes Iran the core of a very grave threat to the U.S. global empire which I believe has already developed far enough to prevent the U.S. from continuing its imperialist suppression of muslim nations with totalitarian, puppet regimes and the use of those to exploit this vast region.
But I think these would be temporary set-backs for Iran and its great alliance of nations lining up against the U.S. Empire. I think the world sees the Giant is crippled severely by its losing wars, by the decline of the U.S. economy relative to the rising nations around the world, the so-called BRICKS, though one must add another I to that for Iran [in addition to India] – BRICKS, and The Great Global Financial Crisis the U.S. produced and is making worse and worse with vast help from its puppets in the EU and elsewhere, as even the “liberal” Brookings Institute just reported on the basis of global indices they analyze.
by Rania Khalek on March 20, 2013
The United States may be finished dropping bombs on Iraq, but Iraqi bodies will be dealing with the consequences for generations to come in the form of birth defects, mysterious illnesses and skyrocketing cancer rates.
Al Jazeera’s Dahr Jamail reports that contamination from U.S. weapons, particularly Depleted Uranium (DU) munitions, has led to an Iraqi health crisis of epic proportions. “[C]hildren being born with two heads, children born with only one eye, multiple tumours, disfiguring facial and body deformities, and complex nervous system problems,” are just some of the congenital birth defects being linked to military-related pollution.
In certain Iraqi cities, the health consequences are significantly worse than those seen in the aftermath of the atomic bombing of Japan at the end of WWII.
The highest rates are in the city of Fallujah, which underwent two massive US bombing campaigns in 2004. Though the U.S. initially denied it, officials later admitted using white phosphorous. In addition, U.S. and British forces unleashed an estimated 2,000 tons of depleted uranium ammunitions in populated Iraqi cities in 2003.
DU, a chemically toxic heavy metal produced in nuclear waste, is used in weapons due to its ability to pierce through armor. That’s why the US and UK were among a handful of nations (France and Israel) who in December refused to sign an international agreement to limit its use, insisting DU is not harmful, science be damned. Meanwhile, the Pentagon’s refusal to release details about where DU munitions were fired has made it difficult to clean up.
Today, 14.7 percent of Fallujah’s babies are born with a birth defect, 14 times the documented rate in Hiroshima and Nagasaki. Fallujah’s babies have also experienced heart defects 13 times the European rate and nervous system defects 33 times that of Europe. That comes on top of a 12-fold rise in childhood cancer rates since 2004. Furthermore, the male-to-female birth ratio is now 86 boys for every 100 girls, indicating genetic damage that affects males more than females.
(On a side note, these pictures are rather sanitized compared to other even more difficult to look at images. See here if you can bear it.)
If Fallujah is the Iraqi Hiroshima, then Basra is its Nagasaki.
According to a study published in the Bulletin of Environmental Contamination and Toxicology, a professional journal based in the southwestern German city of Heidelberg, there was a sevenfold increase in the number of birth defects in Basra between 1994 and 2003.
According to the Heidelberg study, the concentration of lead in the milk teeth of sick children from Basra was almost three times as high as comparable values in areas where there was no fighting.
In addition, never before has such a high rate of neural tube defects (“open back”) been recorded in babies as in Basra, and the rate continues to rise. According to the study, the number of hydrocephalus (“water on the brain”) cases among new-borns is six times as high in Basra as it is in the United States.
This isn’t isolated to Fallujah and Basra. The overall Iraqi cancer rate has also skyrocketed:
Official Iraqi government statistics show that, prior to the outbreak of the First Gulf War in 1991, the rate of cancer cases in Iraq was 40 out of 100,000 people. By 1995, it had increased to 800 out of 100,000 people, and, by 2005, it had doubled to at least 1,600 out of 100,000 people. Current estimates show the increasing trend continuing.
As Grist’s Susie Cagle points out, “That’s potentially a more than 4,000 percent increase in the cancer rate, making it more than 500 percent higher than the cancer rate in the U.S.“
Dr. Mozghan Savabieasfahani, an environmental toxicologist based in Ann Arbor, Michigan, told Jamail that “These observations collectively suggest an extraordinary public health emergency in Iraq. Such a crisis requires urgent multifaceted international action to prevent further damage to public health.”
Instead, the international community, including the nation most responsible for the health crisis (hint: it starts with a “U” and ends with an “S”), is mostly ignoring the problem.
To make matters worse, Iraq’s healthcare system, which was once the envy of the region, is virtually nonexistent due to the mass exodus of Iraq’s medical doctors since 2003. According to recent estimates, there are currently fewer than 100 psychiatrists and 20,0000 physicians serving a population of 31 million Iraqis.
Dahr Jamail was on Democracy Now this morning discussing the horrific effects of military-related pollution in Iraq:
Yanar Mohammad, President of the Organization for Women’s Freedom in Iraq was also on Democracy Now and addressed the toxic legacy of birth defects in Iraq. (I interviewed Mohammed for a piece I wrote for Muftah about the deterioration of Iraqi women’s rights since the invasion, which you can read here.)
If they had a scintilla of decency, Tony Blair, Alastair Campbell and John Scarlett would not show their faces in public again
Ten years ago today, American and British tanks stormed across the border of Kuwait into Iraq, precipitating a torrent of violence which has since cost more than 100,000 Iraqi and Allied lives.
We, the British people, were told by our leader Tony Blair that the invasion was indispensable to Britain’s national security because Saddam Hussein was developing weapons of mass destruction which could be used against us.
Soon after Western forces reached Baghdad, it became plain that no such weapons existed.
Moreover, it also emerged that the Prime Minister had assured President George Bush of Britain’s armed support in deposing Saddam Hussein, well ahead of the WMD claims, because he wished to assist the Americans in doing what he considered a good deed in the world.
Blair’s trusted henchman, Alastair Campbell, colluded with John Scarlett, chairman of Britain’s Joint Intelligence Committee, to produce a document which afterwards proved to be a mass of falsehoods, offering ‘evidence’ of the case for war.
|Alastair Campbell (left), colluded with John Scarlett,
chairman of Britain’s Joint Intelligence Committee, to produce a
document which has proved to be a mass of falsehoods,
offering ‘evidence’ of the case for war
All three men thus committed what seems to some of us a heinous political crime. They concocted a false manifesto to justify taking Britain to war, with the loss of 179 British servicemen’s lives.
Yet a decade on, not only are those responsible walking the streets of London as free men, but they are without shame.
Blair said this week on BBC’s Newsnight that he does not regret the war. ‘If we hadn’t removed Saddam from power, just think, for example, what would be happening if these Arab revolutions were continuing and Saddam, who’s probably 20 times as bad as Assad in Syria, was trying to suppress an uprising in Iraq.’
|Blair, Campbell and Scarlett thus committed what seems to some of us a heinous political crime. They concocted a false manifesto to justify taking Britain to war, with the loss of 179 British servicemen’s lives|
I feel passionate because I was among those duped by the WMD claims.
Before 2003 I wrote many times in these pages that Britain should have nothing to do with a recklessly irresponsible American Republican adventure in Iraq. But then I read the government’s report on Saddam’s weapons. I felt that this had to be believed, and an invasion reluctantly supported.
My wife Penny, who never swallowed Blair and Campbell’s claims, argued bitterly with me. I said pompously: ‘It’s impossible that the Government and the Secret Intelligence Service would lie to us about something this big.’
I was as wrong as I could be. Blair, Campbell and Scarlett made fools of many of us.
What seems to make it all much worse is that they got away with it.
Scarlett, whom Campbell described to the Chilcot Inquiry into the Iraq War as ‘a mate’, was discredited as an intelligence officer by his role in the scandal. He was the professional spook who processed the material he and Downing Street then served up to the nation. He formed an unholy partnership with his political masters.
When it was all over, however, instead of posting him to Ulan Bator as he deserved, Blair rewarded Scarlett’s misconduct by making him director of the Secret Intelligence Service.
Much of the secret service was thoroughly unhappy about the appointment. But the Prime Minister used his power, and Scarlett served his term. Not only that, but having retired from the SIS, he now serves as an ‘independent director’ of The Times newspaper, a role in which Vladimir Putin would be more credible.
Campbell, meanwhile, has become the darling of the BBC, forever a guest on its chat shows and invited to air his views on news programmes as if he was an elder statesman rather than spinmaster of the most mendacious government of modern times.
Campbell’s off-camera behaviour, as a foul-mouthed bully, was brilliantly captured in the political satire The Thick Of It. But the man himself is nowadays welcomed into studios as if he was a national treasure.
Why have we lost our capacity for anger? I suspect that much of the public is content to forgive these people not from Christian charity, but because it does not care much about anything any more.
I spoke at a dinner last week at which a guest asked me how I would behave if Tony Blair walked in. There was obvious surprise and even disbelief when I said that I would decline to shake his hand.
But I meant what I said. It seems outrageous the former Prime Minister is travelling the world, collecting tens of millions of pounds in consultancy and speaking fees, masquerading black-comically as a Middle East ‘peace envoy’, while in Iraq the killing goes on and on. Yet Blair still insists the invasion was a splendid idea.
Amazingly, David Cameron has invited him into Downing Street to give advice — no doubt about how to persuade the British people that we should get stuck into Syria.
Most of Blair’s former Cabinet colleagues, by contrast, have distanced themselves to some degree from the war. John Prescott, Deputy Prime Minister in 2003, said this month that the Iraq invasion ‘cannot be justified’ and offered a bizarre excuse for his support at the time — that he thought George Bush had a plan to resolve the Israeli-Palestinian conflict.
Prescott was, of course, the Blair government’s resident buffoon rather than a serious politician, but at least he has brains enough now to see that he was complicit in a disaster.
Some people think that public and political attitudes to Blair and Campbell (nobody outside Whitehall remembers Scarlett’s name, luckily for him) will change when the interminable Chilcot Inquiry’s report is published later this year.
I will believe that only when it happens. Official inquiries have a long history of making their conclusions so equivocal that, somehow, nobody ends up getting blamed.
A dispute is continuing between Chilcot and the Cabinet Secretary, who has refused to release pre-war correspondence between Blair and President Bush deemed likely to prove that the Prime Minister signed up for the Republicans’ war long before the issue of WMDs was even raised.
Whatever the outcome of this argument, I shall be amazed if Chilcot’s report displays the courage and clarity of thought to say what most of the British people know already — that Blair and his associates are guilty as hell.
Some people say: ‘Oh, but it’s all a long time ago and, anyway, we are out of Iraq now.’
True, but the nightmare Bush and Blair created continues, with scores dead in a bomb blast in Baghdad only yesterday.
The two Western ‘crusaders for freedom’ achieved the remarkable feat of precipitating the killing of almost as many Iraqis as Saddam Hussein accounted for during his murderous rule.
Yet Tony Blair now lives in a splendour that seems to satisfy even his wife’s sybaritic tastes. Alastair Campbell is said to be advising Ed Miliband on how to get into Downing Street and will probably end up with a peerage for his pains.
The former prime minister and his spin doctor have wrought such tragedy and grief in the world that they should be regarded as pariahs. But people choose to forget.
Political rage focuses instead on phone-message hackers from the Press. Yet at least their crimes, repulsive as they were, never killed anyone. That is more than can be said for those of the former tenant of Downing Street.
"Now, in 2013 another band of manipulators is seeking to take the US to war in Syria. Americans – Beware!" |
I. E. Coop
In such a wide field as pastures and crops it is impossible in the time and space available, to do more than summarise the main lines of research and development of recent years having special reference to underdeveloped countries. In the long-term the greatest advances have been made in fundamental knowledge, such as - the basic physiology and biochemistry of plant growth, the significance of the C4 pathway in tropical plants, genetic engineering, embryo transfer in animals, remote sensing photography, the ecology of the world's grasslands. The pressing problems of rangeland degeneration, and social and economic change in human societies, are also better understood.
The immediate problem and the task of this paper is to descend to a lower practical level of what has been learned about the possibilities of increasing pasture and forage production, and of utilising feed grown in the most efficient manner.
Approximately one half of the sheep and three quarters of the goats in developing countries are within the tropics, the remainder being in a band from North Africa through the Near East to China. It is useful to have some classification of climatic zones governing plant growth, and for the tropics it is given below:
|Zone||Rainfall1 (mm)||Rainfall2 (mm)||Growing Period (days)||Dry Season (months)|
|Arid||< 400||< 500||< 90||> 8|
|Semi-arid||400 – 750||500 – 1000||90 – 180||6 – 8|
|Sub humid||750 – 1200||1000 – 1500||180 – 270||4–6|
|Humid||> 1200||> 1500||> 270||< 4|
1 From Unesco (1979)
2 From Jahnke (1982) for Africa
A feature of research and development work in the tropics, with the exception of the arid zone, is that cattle rather than sheep and goats have been used, so that in the absence of good data on small ruminants one has unfortunately to interpret from cattle data. It is proposed to discuss briefly the recent pasture utilisation studies with sheep in the temperate zone and then move to the arid and semi-arid pastoral zones, followed by the cropping/livestock situation in the semi-arid and subhumid zones and finally to the subhumid and humid zones.
4 RD, Christchurch, New Zealand.
PASTURE UTILISATION IN THE TEMPERATE ZONE
Research activities, and the methods derived therefrom, for increasing ruminant production in developed countries follow fairly standarised lines - breeding and selection for improved cultivars, determination of plant nutrient (fertiliser) requirements, determination of animal feed requirements, grazing management studies aimed at integrating pasture growth and animal requirements with maximum efficiency on a year round basis. In the temperate zone under favourable conditions grass-legume pastures are capable of yielding 10–20t DM/ha/ annum on cultivatable land and 5–10t DM/ha/annum on hill land oversown with clovers especially (Trifolium repens).
While these studies have progressed on a broad front, in recent years special attention has been devoted to the most difficult and complex area - the efficient use of pasture under-grazing (Morley, 1981; Parsons et al., 1983). Techniques have been developed for measuring herbage mass, ratio of green to dead leaf, net DM growth of green material under various grazing pressures, the intake of the grazing sheep (or goat) at various levels of pasture availability, the pasture availability needed to promote given levels of production, and the residual pasture (DM/ha) at which production falls below critical levels. Concurrently research has determined the critical and non-critical nutritional periods in the annual cycle of the ewe, and the extent to which the resilience of the ewe to gain and lose body fat may be used to buffer seasonal peaks and troughs of pasture growth (Coop, 1982; Milligan, 1983).
Finally to put this into practice requires control of pasture growth and control or rationing of intake of the grazing sheep. This can only be achieved by adequate subdivision with fencing. This has been greatly facilitated by the development of electric fencing and in really intensive systems by the additional use of cheap portable electric fencing. Such fencing is also used for strip grazing of forage crops, in order to get maximum utilization of the forage.
The efficient conversion of pasture to animal production is a highly complex matter because of the interactions between the grazing animal and the pasture, interactions which vary with season. At low stocking rates percentage utilisation is low and continuous grazing is as good as, or better than, rotational grazing. However to obtain high animal production per hectare intensive grazing at high stocking rates is required in order to give a high percentage utilisation. In this case rotational grazing is superior. In practice compromises become necessary because pasture growth is seasonal, there are periods when utilisation is sacrificed in the interests of achieving high individual animal growth rates and others, such as in winter, when utilisation is much more important than liveweight gain.
The efficiency of utilisation of native pasture in extensive grazing systems running less than 2 sheep/ha is estimated in recent research to be below 30%. When such pasture is improved by oversowing, fertilisation and fencing, utilisation can be increased to 60–70%. On really intensively grazed cultivated pastures efficiencies of 70–85% are possible on a year-round basis. Some appreciation of intensive grazing and utilisation may be gauged from the current practice of wintering pregnant ewes in New Zealand, where the ewes are rotated, at a density of 1000 ewes per hectare, on a daily shift behind electric fences.
In the Northern Hemisphere where winters are colder, greater reliance on hay and silage is made for winter feed. It is estimated that the percentage utilisation of metabolisable energy (ME) of the original pasture, when consumed as hay, is below 50%. For this, and for reasons of cost, the Southern Hemisphere grazing countries place emphasis on utilisation of pasture by the grazing animal with minimal use of conserved fodder.
Advances in grazing management and utilisation have nevertheless been made in northern countries. One example of this is the “two pasture” system developed for the wet cold hill country of Britain, whereby a smaller area of improved pasture is integrated with the larger area of unimproved hill land and utilised at strategic points in the annual cycle of the ewes. Another is the “three pasture” system aimed at minimising worm parasite problems. Finally in all the major grazing countries there is increasing evidence that cattle, sheep and goats can all be beneficial to one another, the special grazing characteristics of each being complimentary to the other. With coats this is seen especially in their preferences for weeds and pasture species not relished by sheep.
If proof is needed that modern pasture/sheep technology can lead to increased animal production the case of New Zealand may be quoted, where from approximately the same area of grazing land, sheep numbers have increased from 33 millions in 1950 to 53 millions in 1965, to 70 millions in 1982 with proportionate increases in meat and wool output. While only some of this temperate zone technology is immediately transferrable to developing countries, the objectives and the principles certainly are.
ARID AND SEMI-ARID ZONE RANGELANDS
The extensive rangelands of the arid and semi-arid zones of developing countries and the peoples they support are in varying degrees of crisis as a result of rangeland degradation, brought about by overstocking. The area is traditionally used solely by pastoralists under nomadic and transhumant systems, but the pressure of human population has led to the incursion of agriculturalists with their livestock into marginal areas, so putting an unbearable pressure on the rangeland vegetation.
Much has been written about the current state of rangeland vegetation, the social and economic impediments as well as the technical difficulties in reversing the deterioration (e.g. Unesco, 1979; Jahnke, 1982; Harrington, 1982 and Malechek, 1982). While there are cases or instances of potential improvements or improvements actually made, the concensus of opinion of authors is that the only solution short and midterm is to reduce grazing pressure. It is recommended that this be achieved by destocking, or by deferred grazing or some other form of grazing management which would permit a more even grazing and reduce severe overgrazing on critical areas. A recent FAO review (FAO 1984) commented that there is need for rehabilitation by the introduction of good management, that forage cultivation is not yet generally accepted and conservation of hay and silage rarely practised. There is a need to introduce forage trees and browze shrubs, but there was little likelihood of increasing forage availability in the near future due to pressure of livestock combined with the persistence of drought.
The productivity of the arid and semi-arid zone rangelands is low. Jahnke (1982), quoting other authorities, gives a figure of 2.5kg DM/ha/annum per mm rainfall, or It DM/ha/annum at 400 mm which is likely to be inefficiently utilised. Such yields cannot hope to generate enough income to provide incentive to introduce improved species even if this were technologically feasible.
While acceptance by the inhabitants and by Governments that reduction in grazing pressure is the only short term solution, one must not be entirely negative. Observation and development project results indicate that there are avenues for improvement and some specific examples of these are listed below.
Grain yields and sheep production were twice as great in South Australia through replacing fallow with subterranean clover and medic pasture, compared with Algeria having a similar Mediterranean climate but not integrating crop and sheep grazing (Allden, 1982, quoting Carter).
In the Drought Prone Areas Programme in Western India the introduction of Cenchrus ciliaris and Lasiurus sindicus increased DM yield from 0.4 t to 3t/ha/annum (Jain, 1983).
Depleted rangeland in China has been shown to be capable of yielding 3t DM/ha/annum by oversowing with milk vetch and fertiliser (Chinzagco project, pers. comm). In another site having 300 mm rainfall, all in summer, the yields of native grassland have been doubled with fertiliser alone, while in cultivated areas the use of newer cultivars of sorghums, maize, and annual grasses for silage, and native grass for hay has also doubled the number of stock carried as well as improving them greatly (FAO 1983).
The Syrian Arab Republic Rangeland Conservation and Development Project is one of the best known, reviving the ancient “Hema” system of grazing control, introducing Atriplex spp. planting fodder trees and creating lamb fattening cooperatives (Draz 1978).
The wide ranging development project in Morocco where Agropyron elongatum has been introduced into a Stipa-Artemisia ecosystem in a 300 mm rainfall area (El Gharbaoui, 1984).
The introduction of Atriplex and Kochia spp. in Saudia Arabia (Hassan 1984),
The legumes Stylosanthes humilis and to a lesser degree S. guyanensis have been shown to be capable of being oversown or direct drilled on sites in the semi-arid zone.
There are also arid or semi-arid rangelands in the temperate zone (U.S.A., South America, South Africa, Australia) which have also degenerated under overstocking during the last 100 years and it is significant that in all of these stock numbers have declined. The most intensively studied are those in the U.S.A. and in a recent review of rangeland management and reseeding results, it is commented that “a considerable portion of western rangelands currently support vegetation assemblages greatly below their potential” (Herbel, 1984; Young et. a!., 1984). Wilson, A.D. (1982) in another review concludes that “there are no technological improvements in the pipeline that will lead to major productivity gains. The basic restrictions of sparse vegetation, low rainfall and a harsh climate are not subject to technological innovation”. Nevertheless there are instances that in all of these countries improvements are technically possible. To take but one example, Stevens and Villalta (1983) at high altitudes in Peru were able to establish ryegrass-clover pastures and to direct-seed lucerne into rangeland with large increases in sheep numbers carried.
The problem is that research and development projects in both developed and developing countries on which the possibilities if improvement have been shown, have high inputs of technical and economic aid. Whether they can survive in a straight commercial sense and whether it is economic to attempt to increase production is highly dubious. In the more favourable sites it may be so, but for most of it, the problem is to halt further deterioration. The poor income-generating power of the extensive rangelands dictates that any improvements must be ecologically sound and low cost, and should act in a catalytic role to permit better utilisation of the much larger area of unimproved land.
Research priorities suggested should include grazing management studies to provide more even grazing pressure, forage conservation, selection of species and cultivars extending growth into the dry period, integration with cropping systems. (Unesco, 1979; Malechek, 1982; Butterworth et.al,1984).
CROP - LIVESTOCK INTEGRATION
Crop production is an occupation of agriculturalists living in villages mostly in the semi-arid and subhumid zones. Traditionally some nomads have included the grazing of crop stubbles in their annual movement, while transhumant pastoralists have also made use of stubbles and crop residues during the dry period. The increasing sedentarisation or semi-sedentarisation of nomads and transhumants, together with movement of agriculturalists with their own livestock in the opposite direction into drier areas, is reducing the areas available for grazing and also increasing the risks of crop failure. The integration of cropping with sheep and goats is primarily in the semi-arid zone but extends into the subhumid zone. Although the cropping regime yields more DM/ha in the form of stubbles, straws and byproducts available for stock the increases in stock numbers more than offsets this. Nevertheless cropping systems and the more intensive and settled human existence in villages or permanent abodes, offers an environment much more amenable to technological change and improvement than does the rangeland. The following research developments in recent years are some of the more promising.
The breeding of improved cultivars of human feed crops - wheat, maize, sorghum, groundnuts etc. and research on fertiliser responses, together with an appreciation that in subsistence agriculture, fertilizer put on crops increases yield sufficiently to release land for planting in animal forage crops.
Research and demonstration has shown that forage production can be expanded considerably by inter-row sowing of legumes with the cereal, using improved cultivars of forage species, and especially replacing the traditional fallow with sown perennial or annual forage crops. Legumes such as Stylosanthes and vetches, and other tropical legumes in higher rainfall areas, are much preferred since their nitrogen level and nutritive value are high and they increase soil nitrogen for the next cereal crop. High yields have been obtained in Cyprus from barley and barley/vetch forage made into hay (Osman and Nersoyan, 1984; Unesco, 1979; FAO, 1983). If a move to greater use of forage crops and more efficient use of grazing stubbles is to be made then control of the sheep and goats becomes important. Attempts should therefore be made to gain acceptance of the electric fence by herders and cultivators.
Intensive fattening of lambs and kids, on locally grown roughage plus concentrates and byproducts, has a double advantage of controlled marketing with a superior product and more importantly of removing young animals to be fattened from the overgrazed rangeland, thereby reducing the grazing pressure. Lamb fattening trials have been reported from several countries showing typically that weaned lambs make gains of 100 – 250 g/day with feed conversion ratios of 6 to 10 according to the energy content of the diet. There is a need to examine what effect this has on the total system.
Some arid and semi-arid areas have water available for irrigation, which is used mainly for cereal or cash crops (cotton) but some is available for forage. Water from the Nile is used in Egypt and Sudan, underground water in Libya and Saudi Arabia. Extremely high yields of lucerne (Medicago sativa ) and Egyptian clover (Trifolium alexandrinum) are obtained and provide a high protein source for cattle, sheep and goats.
Improving the utilisation of low quality roughages is also possible. Low protein levels characteristic of tropical forages during the long dry period are a limiting factor in animal intake and performance. (Minson, 1982). A considerable amount of research work has been done over the last 20 – 25 years on the use of urea to improve the voluntary intake of straws and other low quality roughages by cattle, sheep and goats. Trials conducted in pens have almost universally given good results but selective grazing by animals in the field has caused some doubts about its application in a grazing context (Coombe 1981). A more recent discovery is that alkali or ammonia treatment of straw can increase digestibility by 10–15units, e.g. from 45% to 55–60%. Encouraging results are being obtained from the technique at both the village level (Dolberg et.al., 1981), and the factory level (Creek et.al., 1984).
A much better understanding of protein requirements of sheep and goats has been developed during the last decade, with recognition of the significance of rumen nondegradable protein. This is of special importance in the tropics (Lindsay, 1984).
The outlook then for improvements in pasture and crop production, and of utilisation by sheep and goats in the cropping areas is reasonably encouraging. Whether it can keep pace with the increases in human population is another matter. Fortunately much of the research done in developed countries is less sensitive to environment in a cropping activity than in a grazing activity, and is therefore more likely to find application in the cropping scene. The most important fields of research in the cropping areas as far as sheep and goats are concerned are likely to be further integration of pastoral ism with cropping, conservation and forage production for the dry period, and improvements in the utilisation of straws.
Somewhat similar problems exist in the semi-arid/cold regions of the world such as in the arc from Turkey to China. Here the winter replaces the dry period of the tropics. In the USSR and Northern China for example, many pastoralists have been semi-or wholly sedenterised, and winter bases exist in villages or have been especially constructed. The growing of forage, partly for grazing but mostly for conservation as hay and silage, is a dominant feature of the system (Demiruren, 1982).
SUBHUMID AND HUMID ZONES
Though the line of demarkation between the semi-arid and subhumid zones is diffuse, there is a distinct trend towards tree crop agriculture as well as cropping, towards tall-grass pasture species and a greater density of villages, especially where associated with rice culture. This is accompanied by a shift in the relative importance of large and small ruminants. Whereas in developing countries sheep and goats outnumber cattle by nearly 2 : 1 in the arid and semi-arid zones, cattle outnumber sheep and goats in the subhumid and humid zones. As far as sheep and goats are concerned there are no longer any pastoralists and nearly all the animals are associated with village and cropping agriculture.
Tropical Pasture Development
Present native pastures consisting of Hyperrhenia, Andropogon, Themeda and many other species exist in a savanna landscape derived from forest or woodland. Soils are heavily leached, grazing is primarily with cattle and fire plays an important part in the grass, scrub, tree balance. The most important development in this area in the last few decades has undoubtedly been the selection, breeding and cultivation of improved cultivars of tropical grasses and legumes. The legume is particularly important because of the low nitrogen status of tropical soils. Though this work has been carried out in several tropical environments the driving force has been the CSIRO Division of Tropical Pastures in Queensland, Australia (Mannetje, 1982; Minson, 1982). Now there are established cattle ranches and cattle projects in most tropical countries with rainfall in excess of 800–1000 mm.
Unfortunately, in relation to sheep and goats, the basic grazing experiments and present projects are almost wholly involved with cattle. There are good reasons for this cattle dominance, but not for the exclusion of small ruminants. Very high yields of pasture DM are attainable - up to 30 – 40t/ha/annum but control of pasture growth, maintenance of the grass-legume balance, and ingress of weeds do present greater problems than with temperate pastures (Mannetje, 1982). Nevertheless the potential of these tropical pasture species for small ruminants with or without cattle should be explored. Some trials using sheep and goats have been recorded (Boulton and Norton, 1982; Potts and Humphreys, 1983; Susetyo et.al. 1983) but not yet on a farm scale. Some of the improved species, especially legumes such as Stylosanthes humilis and S. guyanensis, Macroptilium, Desmodium spp are also finding use as forages for establishment on fallows which are grazed by sheep and goats in both semi- and subhumid zones.
Sheep and Goats in the Village
The place of sheep and goats in the village is much the same as described for the arid and semi-arid zones. The scope for increases in numbers and production especially of goats has been emphasised (Devendra, 1980; Roy-Smith, 1982; Zulkifli et.al., 1980; Wilson, R. T. 1982), involving greater utilisation of the considerable byproducts available from cereal and tree crop production, the introduction of improved grass and legume species in available land and recognition of the part which the milking goat rather than the cow could play. The opportunity to exploit the production of tropical grasses and legumes is facilitated by the fact that many sheep and goats, especially in the humid zone, are fed under a cut and carry system, or are let out during the day time on a controlled grazing system. There is no doubt that small ruminant production in this zone could be increased and there are good reasons why it should. Most of the arguments about the relative merits of cattle, sheep and goats are based on personal factors, and on calculations of what ought to be, and there is a need for comparisons to be made on a strictly scientific basis. It is encouraging that recognition of animal production potential in the humid tropics is evidenced by the creation of the joint Australian/Indonesian Centre for animal research and development in Ciawi, Indonesia in 1975.
Utilisation of Pasture in Plantations
In the humid tropics there are large areas of tree crops such as coconut, rubber and oil palm. They are established in association with a tropical legume cover crop which in time regresses to grasses and weeds. Except in coconut plantations often grazed by cattle the herbage available is generally not used at all. Attention has been given by the Rubber Research Institute of Malaysia (Tan and Abraham, 1980) to using sheep to consume this herbage and to reduce the high cost of weed control. Promising results are being achieved, confirming that a considerable potential exists for the utilisation of this large feed resource.
Forage Trees and Tree Byproducts
The utilisation of edible trees and shrubs, and of tree byproducts such as leaves, pods and seeds has received considerable attention in recent years. The characteristics and feed values of tree crops have been reviewed recently by Hutagalung (1981). Particular attention has been given to leguminous trees-such as Leucaena, Gliricidia, Tagasaste spp. since the leaves of leguminous trees, and especially L. Leucocephala, have protein levels in excess of 20%. Recent evidence (Bamualim et.al., 1984) shows that the protein is a good by-pass protein, capable of enhancing intake of low quality roughages where these form the main diet. Acacia spp prevalent in many parts of the tropics can also be valuable in droughts (Snook, 1984).
Leguminous trees and edible shrubs are not confined to the subhumid and humid zones, but can contribute also in the semi-arid zone, and are likely to be utilised much more than in the past.
There is considerable scientific evidence that both pasture and animal production can be increased substantially in all except the arid zone. The major problems are lack of capital to implement improvement and uncertainty about its viability, economically. The greatest and quickest responses in pasture and forage production are likely to be in the subhumid zone, and also in cropping systems.
To maintain impetus in research it is suggested that priority be given to:
breeding and selection of pasture species for the semi-arid zone aimed at increasing the length of the growing period. Selection of low phosphate demanding legumes for this zone.
studies of the utilisation of tropical pastures by sheep and goats.
pasture conservation for the dry period, forage production on cropping land.
studies to increase the complementarity (or integration) of rangeland and cropping land.
continuation of studies aimed at improving the utilisation of low quality roughages by sheep and goats.
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Short Watershed Courses
University of British Columbia
Integrated Watershed Management
"Integrated Watershed Management: A Hyper-Media CD-ROM" was developed for use with the University of British Columbia graduate level Internet course, the core course mentioned in the UBC Watershed Certificate Program above. Issues relating to urban, rural, agricultural, forestry, groundwater, and stream water are treated in an interdisciplinary manner. Basic theory is presented, specific land use and watershed issues are integrated, and numerous case studies are described. The CD is reportedly ideal for practicing professionals, graduate students, resource managers and planners and as a source book for watershed partnerships. Containing more than 700 computer frames, 400 images, numerous graphics, text and over 400 searchable references, the CD costs $100 (Canadian).
Contact: Institute for Resources and Environment, 2206 E. Mall, Univ. of British Columbia, 5997 Inonia Dr., Vancouver BC V6T1Z1. Phone (604) 822-1450; fax (604) 822-1499. Web: http://www.ire.ubc.ca
The California Watershed Academy
Teaching Resource Professionals
About Watershed Processes
Pete Cafferata, Forest Hydrologist
Calif. Dept. of Forestry and Fire Protection (CDF)
Jim Steele, Biologist
Calif. Dept. of Fish and Game (CDFG)
The increase in the number and complexity of the California's state forest practices rules in the last ten years has been dramatic. Hundreds of changes, additions and many new concepts have added to the complexity of designing, approving and conducting a timber harvest operation. Pushing this increase is the Endangered Species Act (ESA) protection extended to the northern spotted owl, marbled murrelet, red legged frog, coho salmon and steelhead trout, and controversy over old growth forest logging. A forester today is expected to know more and have a technical background far greater than expected a few years ago. Often, natural resource protection concepts embedded in a rule are not understood by either the forester, decision maker or the public. A way of advancing new concepts and pragmatic methods was needed.
Five versions of the "Watershed Academy" have been presented to about 200 resource professionals over the past four years. Primary goals have been to provide up-to-date information on key processes affecting aquatic habitats, and to improve students' hydrologic IQ and their analysis and risk assessment skills. Resource professionals completing the academy are trained to understand the basics of natural resource protection law theory, fluvial-geomorphological processes, determine current watershed condition, evaluate whether proposed practices will adversely impact aquatic resources, and develop appropriate mitigation measures to both avoid or minimize additional impacts and accelerate recovery from past practices.
The first three watershed academy sessions were State interagency efforts while the last two were presented by University of California Cooperative Extension (UCCE). Mr. Gary Stacey, CDFG, was the lead person for the interagency efforts managed by Jim Steele and Dr. Richard Harris, UCCE, led the more recent academies. The first sessions using some paid instructors were funded jointly by CDFG and CDF augmented by voluntary instructors from several agencies. The latter sessions with mostly paid instructors were funded by a grant from the National Fish and Wildlife Foundation to the National Marine Fisheries Service (NMFS) and some staff support from CDFG.
The curriculum for the Watershed Academy has been modified for each session based on participant feedback and decisions made by lead persons for each academy session. The lecture portion of the first academy covered the broadest range of topics: fish biology/life history, aquatic amphibian biology/life history, benthic macro-invertebrates, hydrology, fluvial geomorphology, geology, soils, watershed processes, risk assessment for biological and physical parameters, problem synthesis, monitoring, data management and natural resource law. The field portion covered road problem recognition, field assessment techniques for hillslope and instream issues, and stream parameter measurements. Topics covered in the second academy included: natural resource law, fish biology, hydrology, geology, fluvial geomorphology, and watershed assessment procedures. Field topics included road inventory techniques, field examples of differing fish habitat and streamside quality, and demonstrations of instream monitoring techniquesóincluding rapid bio-assessment procedures.
The more recent academies held in the fall of 1998 differed somewhat in philosophy. Rather than emphasizing background information on various watershed related topics, students were asked to study a comprehensive binder of readings to provide a basic level of understanding prior to attending the course. Presenters were primarily from universities, consulting firms and the US Forest Service, with less emphasis on lecturers from state agencies. Lecture topics included: fluvial geomorphology, hydrology, hillslope stability, fish habitat requirements, monitoring, and management issues related to roads and crossings, water quality, and riparian zones. Field exercises stressed hillslope stability assessments, road inventory techniques, and assessing off-site impacts of timber operations.
General impressions from the various academy sessions include the following points: 1) readers and/or binders with large numbers of watershed related papers are taken home for reference, but usually not studied before or during the session itself (however this information was valuable as reference for future timber harvesting plans), 2) the most successful learning experiences take place in the field where practical discussions occur, 3) it is difficult for all the students in a given class to be brought up to a minimum level of understanding for all the background disciplines related to watershed processes (both physical and biological), 4) lecture material should cover half day segments, with the remainder of the day spent in the field illustrating points covered in lecture, 5) no more than one-week should be spent on the academy because of other time commitments for students, and 6) all material presented should be as practically oriented as possible, with specific examples provided to illustrate points covered in lectures. Material presented by Dr. Bill Weaver of Pacific Watershed Associates on issues related to roads, landings, and watercourse crossings has often been cited as the most important information received by the students in each session.
The important concept to be advanced, through understanding watershed function, is how to conduct land management activities compatible with natural resources. Since all natural resource laws are similar in their intent but differ in decision process, having the ability to focus on desired outcomes can streamline a knowledge based permit process, and lower government and private sector costs. Water codes, endangered species protection, and public process can all be included in the same permitting effort if there is confidence in the project's outcome. The Watershed Academy is part of a larger package which attempts to achieve this confidence, including the watershed process courses for professionals and baseline watershed analysis, database consolidation, field studies of watershed principles, and effectiveness monitoring programs.
It is unclear at this time how future offerings of the watershed academy will be presented. Clearly, there is a need for this type of information as most of the North Coast watersheds are listed as impaired water bodies by the U.S. Environmental Protection Agency (EPA) and require development of Total Maximum Daily Load allocations (TMDLs). Additionally, state and federal listings of threatened and endangered species that depend on properly functioning aquatic habitats make this type of training very important. Eventually, CDF would like to have all of its Forest Practice staff participate in the academy. Many private RPFs have also expressed an interest in future sessions. Options will be discussed by state and federal agency managers in 1999.
Contact: Pete Cafferata at (916) 653-9455 or Jim Steele at (916) 653-6194.
The Council of State Governments
Working at a Watershed Level (Interagency Course)
Summary: Oneweek course covering all facets of watershed work, including stream ecology, system dynamics, assessment and analysis, planning methodologies, restoration/management techniques, public involvement strategies and outreach program development.
The Council of State Governments Center for Environment and Safety and other regional organizations are now offering a new watershed training curriculum. Working at a Watershed Level was developed by a consortium of federal agencies, state/local groups and private organizations to improve crossagency watershed training. The course is designed as an introductory level basic training program for agency personnel newly assigned to watershed teams, veteran watershed managers in need of a refresher course and members of citizens groups interested in a cooperative approach to watershed issues.
The Interagency Watershed Training Cooperative, composed of representatives from the U.S. Environmental Protection Agency, Natural Resources Conservation Service, U.S. Fish and Wildlife Service, Bureau of Land Management, Bureau of Reclamation, U.S. Forest Service, and U.S. Army Corps of Engineers, provided leadership for developing the course outline. The Council of State Governments, International City/ County Management Association, Ecological Society of America and other partners assisted with final curriculum design and content.
Working at a Watershed Level covers the principles of watershed ecology, system dynamics, assessment and analysis, planning methodologies, restoration/management techniques, public involvement strategies and outreach program development. The course provides a basic but very broad foundation for considering both ecological and socioeconomic issues in watershed work across a wide range of public and private organizations. One of the motivating forces for developing the course was the need for a more cooperative, coordinated approach to watershed management and a common orientation to the science and societal issues involved.
While it is recognized that state and federal agencies will continue to have unique needs and somewhat discrete processes for watershed planning, management and restoration, it is hoped that Working at a Watershed Level will help to develop a broad, common framework capable of accommodating the disparate interests that may be involved. Public agencies and private interests can only benefit by working together within a watershed, though each may have slightly different approaches and requirements. Ideally, agencies and other stakeholders will be able to use Working at a Watershed Level to identify areas where multiagency interdisciplinary teams can work together on management issues while retaining the ability to satisfy organizational, statutory or regulatory needs.
An outline of the course can be found at the U.S. EPA's Watershed Academy web site at http://www.epa.gov/OWOW/watershed/wacademy/interfed/shedcors.html
1999 Course Schedule
June in the San Joaquin River Basin in California (tentative).
Early September, 1999 in Lafayette, Indiana (tentative).
Late September, 1999 in New England (tentative).
Contact: Barry Tonning, Environmental Policy Analyst, The Council of State Governments, P.O. Box 11910, Lexington, KY 405781910; (606) 244-8228; fax: (606) 244-8239. E-mail: email@example.com . For updates on the training schedules: http://www.statesnews.org/ecos/
The "Working at a Watershed Level" Course
A Student Review
by William Zigler
CSU Fresno Senior in Geography
I had the opportunity to attend the January 11-15, 1999 session of the Working at a Watershed Level training course held in Chico, California. Working for the Sequoia National Forest as a trainee in hydrology with a background in geography, I had hoped to gain enough foundational information to get me started "out in the field". The course exceeded my expectations, progressing systematically from structural stream elements to geomorphology to public outreach and education strategies. It was a great experience! Because of the course's potential benefit, I would like to provide a synopsis.
Working at a Watershed Level is a watershed management training course conducted on the California State University Chico campus, through which flows Big Chico Creek. This course received much popular support in Northern and Central California, drawing many from local, state and federal agencies as well as from consulting firms, local watershed conservancies, irrigation districts, and agriculture. Not all attendees were local: some came from the distant points of West Virginia, Vermont, and Vancouver, British Columbia. In fact, the course was so popular that organizers were forced to limit seating to just over 100 while maintaining a waiting list of about 40. (Organizers are considering providing another training session later this year for those who were unable to attend.) The interest in the course reflects an increased awareness of the importance of watersheds today, especially with a growing population and its requisite demands on the environment.
The course was comprised of classroom training sessions, group discussions, field trips, and two evening social gatherings. During a dinner sponsored by The Nature Conservancy, the speaker, Donald Outen of the Baltimore County, Maryland Dept. of Environmental Protection and Resource Management, gave an in-depth presentation on Baltimore County's Integrated Watershed Management Program. He addressed federal nonpoint source pollution control mandates, State initiatives for restoration of Chesapeake Bay, and local priorities for cooperative water quality projects. Mr. Outen was a knowledgeable speaker who demonstrated effective methods of managing an important watershed surrounding a growing urban center.
Other speakers involved with the course were equally credible. They represented a distinguished and diverse group from academia (CSU Chico, Univ. of Montana, Univ. of Washington, UC Berkeley, and Shasta Community College), state and federal agencies (USFWS, Council of State Governments), private industry (Tetra Tech, Inc.), and conservation groups (Center for Natural Lands Management, Tuolumne River Preservation Trust, Mill Creek Watershed, and Stanislaus River Project). The curriculum was effectively organized, progressing logically with each new theme building upon the last. The curriculum provided something for everyone, focusing on the many aspects of watershed management.
While the classroom sessions were informative, I especially enjoyed our training outside the classroom where we observed watershed management practices in action (with varying degrees of effectiveness). Two afternoons were dedicated to field trips to Butte, Big Chico, and Sycamore Creeks. Professors Matt Kondolf, Paul Maslin, and Morgan Hannaford led groups through a watershed assessment on Butte Creek. Butte Creek was an interesting study in the instream effects of flood control measures. The flood control channel was hastily excavated after a major flood event to divert high flows of water downstream during future flood events. The channel was effective as diverting flood waters away from the narrow main channel; however, it created a problem of sediment deposition at the point bar of the main channel, effectively closing the main channel after a flood event. Closing the main channel posed several problems, such as increased stream velocity, altered erosion patterns, and the inability to provide water at the irrigation diversion dam for agriculture. Currently, the main channel has been reopened but repeated manual removal of sediment at the point bar is required to maintain the main channel.
A different field trip took us to a flood control diversion on Big Chico Creek and the flood outflow area on Sycamore Creek. A similar approach to the Butte Creek flood control diversion was used on Big Chico Creekwith similar results. Point bar sedimentation is a problem on Big Chico Creek while Sycamore Creek has experienced significant streambed degradation. The accumulation of sedimentation on Big Chico Creek is manageable, but the control of erosion along Sycamore Creek is less certain. Both field trips provided excellent examples of good intentions (the protection of property from flood impacts) gone bad. When flood control efforts are undertaken without proper scientific research, the results can be more devastating that the flood event itself.
In summary, Working at a Watershed Level was an excellent experience for a novice like myself. I now have a better feel for how the many subsystems within a watershed comprise the whole, and how we are all impacted by the general health of our watersheds. I enjoyed the enthusiasm and professionalism displayed by our host of speakers and appreciate the material produced by those who put their energy and talent into developing the course. I especially valued the opportunity to view first-hand the effects of inappropriate flood control measures. I would heartily recommend this course to anyone desiring a deeper appreciation and understanding of watersheds.
U.S. Environmental Protection Agency
Within EPA's Office of Wetland, Oceans and Watersheds (OWOW) is the Watershed Academy. The Academy offers several of its own training courses supporting watershed approaches and publicizes watershed courses sponsored by others. In addition, EPA's Office of Water sponsors its own courses that are related to watershed management. The Academy courses are not offered every year and they are apparently not always co-taught or co-sponsored through a university or college. Those that are generally available are listed below:
- Watersheds 101: Applied Watershed Management
- Watersheds 102: Statewide Approach to Watershed Management
- Watersheds 103: TMDL Training for State Practitioners
- Watersheds 104: Executive Overview of the Watershed Approach
- Watersheds 105: Watershed Management Tools Primer
- Watersheds 106: Watershed Partnership Seminar
- Watersheds 107: Using Internet Resources
Watershed Academy 2000 Distance Learning
As part of its Watershed Academy, EPA is developing Academy 2000 Distance Learning, using the Internet as a classroom. Academy 2000 consists of training modules on watershed science, effective communications, and organizational management and development. Modules include:
- Principles of Watershed Management
- Watershed Restoration
- Economics of Sustainability
- Monitoring Consortiums
- Watershed Modeling
- Executive Overview of the Watershed Approach
The website contains 13 modules, with at least seven more planned. Running time for each module is about two hours. The website is www:epa.gov/owow/watershed/wacademy/acad2000/index.html
Inventory of Watershed Training Courses
EPA's Watershed Academy has recently produced an Inventory of Watershed Training Courses, updating its earlier document titled "Watershed Academy Catalogue of Watershed Training Opportunities" (May 1997). Listed as a key action in the Clean Water Action Plan, the inventory provides information on available watershed-related training courses sponsored by governmental and non-governmental organizations. The target audience includes federal, state, and local agency staff; tribes, and watershed groups. EPA took the lead in updating this inventory, while working with the Interagency Watershed Training Cooperative, Natural Resources Training Council, EPA's Office of Water Watershed Training Work Group, and others.
Included courses focus on protecting or restoring watersheds (or aquifers), cover important tools used in watershed protection, or address one aspect of the watershed management cycle (e.g., planning, implementation, evaluation). In addition to agency courses, University training courses are also listed. Check the website for its availability: www:epa.gov/owow/watershed/wacademy/catalog.html
INTERDISCIPLINARY WATERSHED CENTERS
University of California at Davis
Center for Integrated Watershed Science and Management
In the Spring of 1998, UC Davis established a center that highlights and expands the campus role in watershed-related research, teaching, and outreach. Under the auspices of the John Muir Institute for the Environment and with guidance from the Commission on the Environment, the new Center for Integrated Watershed Science and Management has a threefold mission:
· develop, coordinate and track current and future campus efforts at integrated watershed study, with an emphasis on expanding the research support base and promoting the leadership role of the campus in watershed issues;
· support existing and future graduate and undergraduate programs that offer integrative watershed education;
· provide knowledge-based services and support to watershed stakeholders and decision-makers and promote the development of university/agency/foundation partnerships in watershed research.
A broad array of public and private institutions have noted that the next frontier in watershed study lies in the integration between disciplines. Authors have noted that failed attempts to manage multiple use in watersheds stem from traditional single agency, single issue approaches. This problem is exacerbated by academic and government institutions who continue to promote hierarchical, discipline-specific research programs in watershed. To date, no academic program has been singled out as successfully promoting such integration and taking a leadership role in this emerging field.
The University of California, Davis, with its diverse intellectual and technical resources, its numerous watershed-related research and teaching programs, its extensive public outreach program, and its long history of collaboration with federal, state and local watershed agencies, believes that it is ideally suited to become a national and international leader in integrated watershed studies.
The Center is intended to be an administrative and intellectual link between the numerous watershed-related programs currently operating on the UC Davis campus as well as watershed decision-makers and stakeholders. The mission of the Center is guided by a Steering Committee, made up of directors or heads of campus programs, and an Advisory Board, composed of watershed experts outside of the UC system. A goal of the Center will be to become self-sustaining within three years through recharge for knowledge-based services, indirect cost return monies, and agency or foundation support.
New faculty positions are being created and filled to expand watershed education and research at UCD. For example, an Assistant Professor in Watershed Hydrology was recently added to the Hydrology Program to develop a quantitative field-experimental research program in Watershed Hydrology. The appointee is expected to lead a team-taught field course in hydrology, an undergraduate course in watershed hydrology, and a graduate level course in experimental watershed hydrology.
The measure of success of the education and teaching mission of the program will be the changes in the curriculum of existing programs on campus, the development of new programs, and the number of students who attend courses associated with the center. An additional measure will be the number of student internships and fellowships sponsored through the center.
Contact: Dr. Jeff Mount, Dept. of Geology, UC Davis, E-mail: firstname.lastname@example.org
Also: Hydrology Program, Dept. of Land, Air and Water Resources, http://lawr.ucdavis.edu
University of Idaho
Eco-hydraulics Research Group
More recent approaches to river management are multi-objective, balancing beneficial uses for power generation, water supply and agriculture with the protection and enhancement of the riverine habitat, water quality, recreational use and aesthetics. These restoration and enhancement approaches place an emphasis on allowing the physical processes to drive the ecological healing by natural evolution, rather than an instantaneous engineering fix. Implementing this restoration philosophy, developing management plans, simulating the hydrological or ecological responses and untangling the complexities of aquatic systems require an interdisciplinary approach, which crosses the boundaries of science and engineering programs. The term "Eco-hydraulics" comes from the new forum created in 1996 by the International Association for Hydraulic Research (IAHR).
To address the specific problems of the Pacific Northwest, the University of Idaho's Eco-hydraulics Research Group comprises faculty from the Departments of Fish and Wildlife Biology, Geography, and the College of Engineering. Faculty expertise includes decision theory, GIS, stream ecology, fisheries biology, biological and microbiological processes, hydrology, hydraulic engineering, sediment transport, geomorphology, computer simulations and computational hydraulics. Collaboration occurs with state and federal agencies, the Idaho Water Resources Research Institute, and international researchers.
Graduate student opportunities exist leading to M.S., M.Engr. or Ph.D. degrees in Civil Engineering, Fish and Wildlife Resources, or Geography. Studies can be undertaken at the Moscow or Boise Campus of the U-I. The Eco-hydraulics laboratory also offers a range of residential short courses for practicing engineers, scientists and planners involved in ecological restoration and natural resource management.
Contact: Eco-hydraulics Research Group, University of Idaho, 800 Park Blvd., Suite 200, Boise ID 83712. (208) 387-1745; Fax (208) 387-1246. Website at: www.engbio.uidaho.edu
Pennsylvania State University
Center for Watershed Stewardship
The Center for Watershed Stewardship is an initiative begun in 1998 that is co-led by the Dept. of Landscape Architecture and School of Forest Resources and funded by a major grant from the Heinz Endowments. Its purpose is "to create the next generation of watershed professionals by combining interdisciplinary capabilities with strong disciplinary bases in a community-oriented context."
Inaugural programs include a graduate option in watershed stewardship and a continuing education program of short courses, seminars, and conferences for natural and water resources professionals and community leaders. Topics to be offered by the Continuing Education and Outreach Program will be: stream corridor management and restoration, community land use planning and design, GIS applications for watershed planning and management, non-profit organizational development and fiscal administration, legal and institutional aspects, water quality management for rural and urban watersheds, natural processes in watersheds, wetland restoration and design, and fisheries restoration and enhancement.
Specific courses will be designed to accommodate a diverse range of participants including landscape architects, community planners, watershed association and land trust staff, community volunteers and activists, natural resource managers, regulatory agency personnel, and environmental consultants, designers, and engineers. Other offerings are being developed through the Center or being co-hosted with other agencies, educational institutions, and departments at Penn State.
Contact: Kerry Wedel, Director, Center for Watershed Stewardship, The Pennsylvania State University, Room 8B, Ferguson Bldg., University Park, PA 16802-4302; (814) 865-8911, fax (814) 865-3725. E-mail: email@example.com
Internet Sites for the Professional Watershed Job Market
American Water Resources Association
Universities Council on Water Resources
COMMUNITY OUTREACH SERVICES
Portland State University
Community Watershed Stewardship Program
By Kristin Schaeffer, Graduate Assistant for CWSP
The Community Watershed Stewardship Program (CWSP) is a partnership between Portland State University (PSU) and the City of Portland's Bureau of Environmental Services (BES) in order to facilitate research and public outreach cultivating stewardship ideology. Our mission is to serve as a catalyst for community ownership of watershed health.
Program goals are to:
- Raise awareness and open avenues for actions to improve watershed health.
- Foster an increased sense of stewardship and promote community initiated projects to improve water quality and habitat in the watersheds of Portland.
- Promote community access to communication systems, training and education about watershed issues and activities.
- Foster partnerships between schools, businesses, interest groups, neighborhood associations, government agencies and individuals in watershed activities.
- Promote citizen evaluation and monitoring of watershed health and activities.
- Strive to be reflective of a diverse community, their values and evaluation.
- Information dissemination and awareness building
- Education and training
- Restoration and enhancement
- Evaluation of the health of watersheds
The role PSU takes in the Stewardship Program is multifaceted. We have two faculty members who are very involved in the management of our program and coordinate graduate assistant (GA) recruitment each year. As well, they facilitate other professor involvement in the Stewardship Program by offering access to our program resources. Since the program works intensively with community efforts, there are watershed research materials and connections we have made with the public. This collective data is then shared with professors who can use the resources to help base their curriculum on.
In the past three years, PSU has instituted a Capstone Class requirement for undergraduate seniors. These classes are intended to get the students involved in community outreach work and offer a number of different classes for specific fields, such as watershed education and awareness. We have found that our program has been getting more involved as a resource for Capstone students and helping to energize stewardship efforts through curriculum development.
Currently, our program funds four graduate assistants who work on average 15 hours a week. Three of the GAs focus on watershed stewardship efforts for a specific watershed, and one GA focuses on the administrative needs of the program and helps to facilitate the Stewardship MiniGrants Program. The GAs come from a variety of backgrounds and departments at PSU, which can change every one to two years. For example, we recruited Joe Blowers who is focusing on a Masters of Science in Teaching Science, Steve Gilchrist is focusing on a Ph.D in the School of Education, Kristin Schaeffer is focusing on a Masters of Public Administration emphasizing on Natural Resources Policy and Administration, and Clint Wertz is focusing on a Masters of Urban Studies and Planning.
At this time, a degree program does not exist in any of the schools or programs at PSU. For the past three years, the Stewardship Program has developed a Watershed 101 course available to communities in which they receive a certificate upon completion. In the past it was a course that was only offered at PSU, but some time was spent last year in developing it for Neighborhood Associations, Watershed Councils and other interested community groups who will have access to the course beyond the campus. We are still in the process of refining this educational component.
For more information, please contact us at phone (503) 8235625, or 'firstname.lastname@example.org'
Oregon State University Extension Service
Watershed Stewardship Guide
OSU Extension's 1998 publication, Watershed Stewardship: A Learning Guide, is intended to help residents and volunteers be good stewards of their watershed. The driving force for the development of this guide was the 1995-97 Oregon Coastal Salmon Restoration Initiative which focused primarily on coho salmon. However, the contents are relevant to all salmonids west of the Cascades. Contained within a two-inch thick 3-ring binder, the material is organized into 3 "user-friendly" sections: IWorking Together to Create Successful Groups, IIUnderstanding and Enhancing Ecosystems, IIIConnecting Resource Management to Watershed Ecosystems.
Copies can be ordered for $32.00 per copy of OSU Extension publication EM 8714 from: Publications Orders, Extension & Station Communications, Oregon State University, 422 Kerr Administration, Corvallis OR 97332-2119. Fax: (541) 737-0817. Discounts are offered on orders of 100 or more copies. Call (541) 737-2513 for price quotes. Also contact them for the latest catalog of publications, software, and videotape programs.
Some other watershed-related materials from Extension include:
- Healthy Watersheds videotape (20 min.) , VTP 019, $20.00
- We All Live Downstream videotape (29 min.) , VTP 021, $30.00
- The Miracle at Bridge Creek Case Study videotape (30 min.), VTP 013, $30.00
- Community Ventures: Interest-Based Problem Solving Process and Techniques, WREP 134, $1.50
- Maintaining Woodland Roads, EC 1139, $1.25
- Water Quality and Our Forests: Western Oregon Research videotape, VTP 014, $25. |
July 2, 2001
ANNALS OF PUBLIC HEALTH
Millions of people owe their lives to Fred Soper.
Why isn't he a hero?
In the late nineteen-thirties, a chemist who worked for the J.R. Geigy company, in Switzerland, began experimenting with an odorless white crystalline powder called dichloro-diphenyl-trichloroethane. The chemist, Paul Müller, wanted to find a way to protect woollens against moths, and his research technique was to coat the inside of a glass box with whatever chemical he was testing, and then fill it with houseflies. To his dismay, the flies seemed unaffected by the new powder. But, in one of those chance decisions on which scientific discovery so often turns, he continued his experiment overnight--and in the morning all the flies were dead. He emptied the box, and put in a fresh batch of flies. By the next morning, they, too, were dead. He added more flies, and then a handful of other insects. They all died. He scrubbed the box with an acetone solvent, and repeated the experiment with a number of closely related compounds that he had been working with. The flies kept dying. Now he was excited: had he come up with a whole line of potent new insecticides? As it turned out, he hadn't. The new candidate chemicals were actually useless. To his amazement, what was killing the flies in the box were scant traces of the first compound, dichloro-diphenyl-trichloroethane--or, as it would come to be known, DDT.
In 1942, Geigy sent a hundred kilograms of the miracle powder to its New York office. The package lay around, undisturbed, until another chemist, Victor Froelicher, happened to translate the extraordinary claims for DDT into English, and then passed on a sample to the Department of Agriculture, which in turn passed it on to its entomology research station, in Orlando, Florida. The Orlando laboratory had been charged by the Army to develop new pesticides, because the military, by this point in the war, was desperate for a better way to protect its troops against insect-borne disease. Typhus--the lethal fever spread by lice--had killed millions of people during and after the First World War and was lurking throughout the war zones. Worse, in almost every theatre of operations, malaria-carrying mosquitoes were causing havoc. As Robert Rice recounted in this magazine almost fifty years ago, the First Marine Division had to be pulled from combat in 1942 and sent to Melbourne to recuperate because, out of seventeen thousand men, ten thousand were incapacitated with malarial headaches, fevers, and chills. Malaria hit eighty-five per cent of the men holding onto Bataan. In fact, at any one time in the early stages of the war, according to General Douglas MacArthur, two-thirds of his troops in the South Pacific were sick with malaria. Unless something was done, MacArthur complained to the malariologist Paul Russell, it was going to be "a long war." Thousands of candidate insecticides were tested at Orlando, and DDT was by far the best.
To gauge a chemical's potential against insects, the Orlando researchers filled a sleeve with lice and a candidate insecticide, slipped the sleeve over a subject's arm, and taped it down at both ends. After twenty-four hours, the dead lice were removed and fresh lice were added. A single application of DDT turned out to kill lice for a month, almost four times longer than the next-best insecticide. As Rice described it, researchers filled twelve beakers with mosquito larvae, and placed descending amounts of DDT in each receptacle--with the last beaker DDT free. The idea was to see how much chemical was needed to kill the mosquitoes. The mosquito larvae in every beaker died. Why? Because just the few specks of chemical that floated through the air and happened to land in the last beaker while the experiment was being set up were enough to kill the mosquitoes. Quickly, a field test was scheduled. Two duck ponds were chosen, several miles apart. One was treated with DDT. One was not. Spraying was done on a day when the wind could not carry the DDT from the treated to the untreated pond. The mosquito larvae in the first pond soon died. But a week later mosquito larvae in the untreated pond also died: when ducks from the first pond visited the second pond, there was enough DDT residue on their feathers to kill mosquitoes there as well.
The new compound was administered to rabbits and cats. Rice tells how human volunteers slathered themselves with it, and sat in vaults for hours, inhaling the fumes. Tests were done to see how best to apply it. "It was put in solution or suspension, depending on what we were trying to do," Geoffrey Jeffery, who worked on DDT at the Tennessee Valley Authority, recalls. "Sometimes we'd use some sort of petroleum-based carrier, even diesel oil, or add water to a paste or concentration and apply it on the wall with a Hudson sprayer." Under conditions of great secrecy, factories were set up, to manufacture the new chemical by the ton. It was rushed to every Allied theatre. In Naples, in 1944, the Army averted a catastrophic typhus epidemic by "dusting" more than a million people with DDT powder. The Army Air Force built DDT "bombs," attaching six-hundred-and-twenty-five-gallon tanks to the underside of the wings of B-25s and C-47s, and began spraying Pacific beachheads in advance of troop arrivals. In Saipan, invading marines were overtaken by dengue, a debilitating fever borne by the Aedes variety of mosquito. Five hundred men were falling sick every day, each incapacitated for four to five weeks. The medical officer called in a DDT air strike that saturated the surrounding twenty-five square miles with nearly nine thousand gallons of five-per-cent DDT solution. The dengue passed. The marines took Saipan.
It is hard to overestimate the impact that DDT's early success had on the world of public health. In the nineteen-forties, there was still malaria in the American South. There was malaria throughout Europe, Asia, and the Caribbean. In India alone, malaria killed eight hundred thousand people a year. When, in 1920, William Gorgas, the man who cleansed the Panama Canal Zone of malaria, fell mortally ill during a trip through England, he was knighted on his deathbed by King George V and given an official state funeral at St. Paul's Cathedral--and this for an American who just happened to be in town when he died. That is what it meant to be a malaria fighter in the first half of the last century. And now there was a chemical--the first successful synthetic pesticide--that seemed to have an almost magical ability to kill mosquitoes. In 1948, Müller won the Nobel Prize for his work with DDT, and over the next twenty years his discovery became the centerpiece of the most ambitious public-health campaign in history.
Today, of course, DDT is a symbol of all that is dangerous about man's attempts to interfere with nature. Rachel Carson, in her landmark 1962 book, "Silent Spring," wrote memorably of the chemical's environmental consequences, how its unusual persistence and toxicity had laid waste to wildlife and aquatic ecosystems. Only two countries--India and China--continue to manufacture the substance, and only a few dozen more still use it. In May, at the Stockholm Convention on Persistent Organic Pollutants, more than ninety countries signed a treaty, placing DDT on a restricted-use list, and asking all those still using the chemical to develop plans for phasing it out entirely. On the eve of its burial, however--and at a time when the threat of insect-borne disease around the world seems to be resurgent--it is worth remembering that people once felt very differently about DDT, and that between the end of the Second World War and the beginning of the nineteen-sixties it was considered not a dangerous pollutant but a lifesaver. The chief proponent of that view was a largely forgotten man named Fred Soper, who ranks as one of the unsung heroes of the twentieth century. With DDT as his weapon, Soper almost saved the world from one of its most lethal afflictions. Had he succeeded, we would not today be writing DDT's obituary. We would view it in the same heroic light as penicillin and the polio vaccine.
Fred Soper was a physically imposing man. He wore a suit, it was said, like a uniform. His hair was swept straight back from his forehead. His eyes were narrow. He had large wire-rimmed glasses, and a fastidiously maintained David Niven mustache. Soper was born in Kansas in 1893, received a doctorate from the Johns Hopkins School of Public Health, and spent the better part of his career working for the Rockefeller Foundation, which in the years before the Second World War--before the establishment of the United Nations and the World Health Organization--functioned as the world's unofficial public-health directorate, using its enormous resources to fight everything from yellow fever in Colombia to hookworm in Thailand.
In those years, malaria warriors fell into one of two camps. The first held that the real enemy was the malaria parasite--the protozoan that mosquitoes pick up from the blood of an infected person and transmit to others. The best way to break the chain of infection, this group argued, was to treat the sick with antimalarial drugs, to kill the protozoan so there was nothing for mosquitoes to transmit. The second camp held, to the contrary, that the mosquito was the real enemy, since people would not get malaria in the first place if there were no mosquitoes around to bite them. Soper belonged to the latter group, and his special contribution was to raise the killing of mosquitoes to an art. Gorgas, Soper's legendary predecessor, said that in order to fight malaria you had to learn to think like a mosquito. Soper disagreed. Fighting malaria, he said, had very little to do with the intricacies of science and biology. The key was learning to think like the men he hired to go door-to-door and stream-to-stream, killing mosquitoes. His method was to apply motivation, discipline, organization, and zeal, in understanding human nature. Fred Soper was the General Patton of entomology.
While working in South America in 1930, Soper had enforced a rigorous protocol for inspecting houses for mosquito infestation, which involved checking cisterns and climbing along roof gutters. (He pushed himself so hard perfecting the system in the field that he lost twenty-seven pounds in three months.) He would map an area to be cleansed of mosquitoes, give each house a number, and then assign each number to a sector. A sector, in turn, would be assigned to an inspector, armed with the crude pesticides then available; the inspector's schedule for each day was planned to the minute, in advance, and his work double-checked by a supervisor. If a supervisor found a mosquito that the inspector had missed, he received a bonus. And if the supervisor found that the inspector had deviated by more than ten minutes from his preassigned schedule the inspector was docked a day's pay. Once, in the state of Rio de Janeiro, a large ammunition dump--the Niterói Arsenal--blew up. Soper, it was said, heard the explosion in his office, checked the location of the arsenal on one of his maps, verified by the master schedule that an inspector was at the dump at the time of the accident, and immediately sent condolences and a check to the widow. The next day, the inspector showed up for work, and Soper fired him on the spot--for being alive. Soper, in one memorable description, "seemed equally capable of browbeating man or mosquito." He did not engage in small talk. In 1973, at Soper's eightieth-birthday party, a former colleague recounted how much weight he had lost working for Soper; another told a story of how Soper looked at him uncomprehendingly when he asked to go home to visit his ailing wife; a third spoke of Soper's betting prowess. "He was very cold and very formal," remembers Andrew Spielman, a senior investigator in tropical disease at the Harvard School of Public Health and the author, with Michael D'Antonio, of the marvellous new book "Mosquito: A Natural History of Our Most Persistent and Deadly Foe." "He always wore a suit and tie. With that thin little mustache and big long upper lip, he scared the hell out of me."
One of Soper's greatest early victories came in Brazil, in the late nineteen-thirties, when he took on a particularly vicious strain of mosquito known as Anopheles gambiae. There are about twenty-five hundred species of mosquito in the world, each with its own habits and idiosyncrasies--some like running water, some like standing water, some bite around the ankles, some bite on the arms, some bite indoors, some bite outdoors--but only mosquitoes of the genus Anopheles are capable of carrying the human malaria parasite. And, of the sixty species of Anopheles that can transmit malaria, gambiae is the variety best adapted to spreading the disease. In California, there is a strain of Anopheles known as freeborni, which is capable of delivering a larger dose of malaria parasite than gambiae ever could. But freeborni is not a good malaria vector, because it prefers animals to people. Gambiae, by contrast, bites humans ninety-five per cent of the time. It has long legs and yellow-and-black spotted wings. It likes to breed in muddy pools of water, even in a water-filled footprint. And, unlike many mosquitoes, it is long-lived, meaning that once it has picked up the malaria parasite it can spread the protozoan to many others. Gambiae gathers in neighborhoods in the evenings, slips into houses at dusk, bites quietly and efficiently during the night, digests its "blood meal" while resting on the walls of the house, and then slips away in the morning. In epidemiology, there is a concept known as the "basic reproduction number," or BRN, which refers to the number of people one person can infect with a contagious disease. The number for H.I.V., which is relatively difficult to transmit, is just above one. For measles, the BRN is between twelve and fourteen. But with a vector like gambiae in the picture the BRN for malaria can be more than a hundred, meaning that just one malarious person can be solely responsible for making a hundred additional people sick. The short answer to the question of why malaria is such an overwhelming problem in Africa is that gambiae is an African mosquito.
In March, 1930, a Rockefeller Foundation entomologist named Raymond Shannon was walking across tidal flats to the Potengi River, in Natal, Brazil, when he noticed, to his astonishment, two thousand gambiae larvae in a pool of water, thousands of miles from their homeland. Less than a kilometre away was a port where French destroyers brought mail across the Atlantic from Africa, and Shannon guessed that the mosquito larvae had come over, fairly recently, aboard one of the mail ships. He notified Soper, who was his boss, and Soper told Brazilian officials to open the dykes damming the tidal flats, because salt water from the ocean would destroy the gambiae breeding spots. The government refused. Over the next few years, there were a number of small yet worrisome outbreaks of malaria, followed by a few years of drought, which kept the problem in check. Then, in 1938, the worst malaria epidemic in the history of the Americas broke out. Gambiae had spread a hundred and fifty miles along the coast and inland, infecting a hundred thousand people and killing as many as twenty thousand. Soper was called in. This was several years before the arrival of DDT, so he brought with him the only tools malariologists had in those years: diesel oil and an arsenic-based mixture called Paris green, both of which were spread on the pools of water where gambiae larvae bred; and pyrethrum, a natural pesticide made from a variety of chrysanthemum, which was used to fumigate buildings. Four thousand men were put at his disposal. He drew maps and divided up his troops. The men wore uniforms, and carried flags to mark where they were working, and they left detailed written records of their actions, to be reviewed later by supervisors. When Soper discovered twelve gambiae in a car leaving an infected area, he set up thirty de-insectization posts along the roads, spraying the interiors of cars and trucks; seven more posts on the rail lines; and defumigation posts at the ports and airports. In Soper's personal notes, now housed at the National Library of Medicine, in Bethesda, there is a cue card, on which is typed a quotation from a veteran of the Rockefeller Foundation's efforts, in the early twentieth century, to eradicate hookworm. "Experience proved that the best way to popularize a movement so foreign to the customs of the people . . . was to prosecute it as though it were the only thing in the universe left undone." It is not hard to imagine the card tacked above Soper's desk in Rio for inspiration: his goal was not merely to cripple the population of gambiae, since that would simply mean that they would return, to kill again. His goal was to eliminate gambiae from every inch of the region of Brazil that they had colonized--an area covering some eighteen thousand square miles. It was an impossible task. Soper did it in twenty-two months.
While DDT was being tested in Orlando, Soper was in North Africa with the United States Typhus Commission, charged with preventing the kind of louse-spread typhus epidemics that were so devastating during the First World War. His tool of choice was a delousing powder called MYL. Lice live in the folds of clothing, and a previous technique had been to treat the clothing after people had disrobed. But that was clearly not feasible in Muslim cities like Cairo and Algiers, nor was it practical for large-scale use. So Soper devised a new technique. He had people tie their garments at the ankles and wrists, and then he put the powder inside a dust gun, of the sort used in gardening, and blew it down the collar, creating a balloon effect. "We were in Algiers, waiting for Patton to get through Sicily," Thomas Aitken, an entomologist who worked with Soper in those years, remembers. "We were dusting people out in the countryside. This particular day, a little old Arab man, only about so high, came along with his donkey and stopped to talk to us. We told him what we were doing, and we dusted him. The next day, he comes by again and says that that had been the first time in his life that he had ever been able to sleep through the night."
In December of 1943, the typhus team was dispatched to Naples, where in the wake of the departing German Army the beginnings of a typhus epidemic had been detected. The rituals of Cairo were repeated, only this time the typhus fighters, instead of relying on MYL (which easily lost its potency), were using DDT. Men with dusters careered through the narrow cobblestoned streets of the town, amid the wreckage of the war, delousing the apartment buildings of typhus victims. Neapolitans were dusted as they came out of the railway stations in the morning, and dusted in the streets, and dusted in the crowded grottoes that served as bomb shelters beneath the city streets. In the first month, more than 1.3 million people were dusted, saving countless lives.
Soper's diary records a growing fascination with this new weapon. July 25, 1943: "Lunch with L.L. Williams and Justin Andrews. L.L. reports that he has ordered 10,000 lbs of Neocid [DDT]and that Barber reports it to be far superior to [Paris Green]for mosquitoes." February 25, 1944: "Knipling visits laboratory. Malaria results [for DDT]ARE FANTASTIC." When Rome fell, in mid-1944, Soper declared that he wanted to test DDT in Sardinia, the most malarious part of Italy. In 1947, he got his wish. He pulled out his old organization charts from Brazil. The island--a rocky, mountainous region the size of New Hampshire, with few roads--was mapped and divided up hierarchically, the smallest unit being the area that could be covered by a sprayer in a week. Thirty-three thousand people were hired. More than two hundred and eighty-six tons of DDT were acquired. Three hundred and thirty-seven thousand buildings were sprayed. The target Anopheles was labranchiae, which flourishes not just in open water but also in the thick weeds that surround the streams and ponds and marshes of Sardinia. Vegetation was cut back, and a hundred thousand acres of swampland were drained. Labranchiae larvae were painstakingly collected and counted and shipped to a central laboratory, where precise records were kept of the status of the target vector. In 1946, before the campaign started, there were seventy-five thousand malaria cases on the island. In 1951, after the campaign finished, there were nine.
"The locals regarded this as the best thing that had ever happened to them," Thomas Aitken says. He had signed on with the Rockefeller Foundation after the war, and was one of the leaders of the Sardinian effort. "The fact that malaria was gone was welcome," he went on. "But also the DDT got rid of the houseflies. Sardinian houses were made of stone. The wires for the lights ran along the walls near the ceiling. And if you looked up at the wires they were black with housefly droppings from over the years. And suddenly the flies disappeared." Five years ago, Aitken says, he was invited back to Sardinia for a celebration to mark the forty-fifth anniversary of malaria's eradication from the island. "There was a big meeting at our hotel. The public was invited, as well as a whole bunch of island and city officials, the mayor of Cagliari, and representatives of the Italian government. We all sat on a dais, at the side of the room, and I gave a speech there, in Italian, and when I finished everybody got up and clapped their hands and was shouting. It was very embarrassing. I started crying. I couldn't help it. Just reminiscing now . . ."
Aitken is a handsome, courtly man of eighty-eight, lean and patrician in appearance. He lives outside New Haven, in an apartment filled with art and furniture from his time in Sardinia. As he thought back to those years, there were tears in his eyes, and at that moment it was possible to appreciate the excitement that gripped malariologists in the wake of the Second World War. The old-school mosquito men called themselves mud-hen malariologists, because they did their job in swamps and ditches and stagnant pools of water. Paris green and pyrethrum were crude insecticides that had to be applied repeatedly; pyrethrum killed only those mosquitoes that happened to be in the room when you were spraying. But here, seemingly, was a clean, pure, perfectly modern weapon. You could spray a tiny amount on a wall, and that single application would kill virtually every mosquito landing on that surface for the next six months. Who needed a standing army of inspectors anymore? Who needed to slog through swamps? This was an age of heroics in medicine. Sabin and Salk were working on polio vaccines with an eye to driving that disease to extinction. Penicillin was brand new, and so effective that epidemiologists were dreaming of an America without venereal disease. The extinction of smallpox, that oldest of scourges, seemed possible. All the things that we find sinister about DDT today--the fact that it killed everything it touched, and kept on killing everything it touched--were precisely what made it so inspiring at the time. "The public-health service didn't pay us a lot," says McWilson Warren, who spent the early part of his career fighting malaria in the Malaysian jungle. "So why were we there? Because there was something so wonderful about being involved with people who thought they were doing something more important than themselves." In the middle of the war, Soper had gone to Egypt, and warned the government that it had an incipient invasion of gambiae. The government ignored him, and the next year the country was hit with an epidemic that left more than a hundred thousand dead. In his diary, Soper wrote of his subsequent trip to Egypt, "In the afternoon to the Palace where Mr. Jacobs presents me to His Majesty King Faruk. The King says that he is sorry to know that measures I suggested last year were not taken at that time." Soper had triumphed over gambiae in Brazil, driven lice from Cairo and Naples, and had a weapon, DDT, that seemed like a gift from God--and now kings were apologizing to him. Soper started to dream big: Why not try to drive malaria from the entire world?
Fred Soper's big idea came to be known as the Global Malaria Eradication Programme. In the early nineteen-fifties, Soper had been instrumental in getting the Brazilian malariologist Marcolino Candau--whom he had hired during the anti-gambiae campaign of the nineteen-thirties--elected as director-general of the World Health Organization, and, in 1955, with Candau's help, Soper pushed through a program calling on all member nations to begin a rigorous assault on any malaria within their borders. Congress was lobbied, and John Kennedy, then a senator, became an enthusiastic backer. Beginning in 1958, the United States government pledged the equivalent of billions in today's dollars for malaria eradication--one of the biggest commitments that a single country has ever made to international health. The appeal of the eradication strategy was its precision. The idea was not to kill every Anopheles mosquito in a given area, as Soper had done with gambiae in Brazil. That was unnecessary. The idea was to use DDT to kill only those mosquitoes which were directly connected to the spread of malaria--only those which had just picked up the malaria parasite from an infected person and were about to fly off and infect someone else. When DDTis used for this purpose, Spielman writes in "Mosquito," "it is applied close to where people sleep, on the inside walls of houses. After biting, the mosquitoes generally fly to the nearest vertical surface and remain standing there for about an hour, anus down, while they drain the water from their gut contents and excrete it in a copious, pink-tinged stream. If the surfaces the mosquitoes repair to are coated by a poison that is soluble in the wax that covers all insects' bodies, the mosquitoes will acquire a lethal dose." Soper pointed out that people who get malaria, and survive, generally clear their bodies of the parasite after three years. If you could use spraying to create a hiatus during which minimal transmission occurred--and during which anyone carrying the parasite had a chance to defeat it--you could potentially eradicate malaria. You could stop spraying and welcome the mosquitoes back, because there would be no more malaria around for them to transmit. Soper was under no illusions about how difficult this task would be. But, according to his calculations, it was technically possible, if he and his team achieved eighty-per-cent coverage--if they sprayed eight out of every ten houses in infected areas.
Beginning in the late fifties, DDT was shipped out by the ton. Training institutes were opened. In India alone, a hundred and fifty thousand people were hired. By 1960, sixty-six nations had signed up. "What we all had was a handheld pressure sprayer of three-gallon capacity," Jesse Hobbs, who helped run the eradication effort in Jamaica in the early sixties, recalls. "Generally, we used a formulation that was water wettable, meaning you had powder you mixed with water. Then you pressurized the tank. The squad chief would usually have notified the household some days before. The instructions were to take the pictures off the wall, pull everything away from the wall. Take the food and eating utensils out of the house. The spray man would spray with an up-and-down movement--at a certain speed, according to a pattern. You started at a certain point and sprayed the walls and ceiling, then went outside to spray the eaves of the roof. A spray man could cover ten to twelve houses a day. You were using about two hundred milligrams per square foot of DDT, which isn't very much, and it was formulated in a way that you could see where you sprayed. When it dried, it left a deposit, like chalk. It had a bit of a chlorine smell. It's not perfume. It's kind of like swimming-pool water. People were told to wait half an hour for the spray to dry, then they could go back." The results were dramatic. In Taiwan, much of the Caribbean, the Balkans, parts of northern Africa, the northern region of Australia, and a large swath of the South Pacific, malaria was eliminated. Sri Lanka saw its cases drop to about a dozen every year. In India, where malaria infected an estimated seventy-five million and killed eight hundred thousand every year, fatalities had dropped to zero by the early sixties. Between 1945 and 1965, DDT saved millions--even tens of millions--of lives around the world, perhaps more than any other man-made drug or chemical before or since.
What DDT could not do, however, was eradicate malaria entirely. How could you effectively spray eighty per cent of homes in the Amazonian jungle, where communities are spread over hundreds of thousands of highly treacherous acres? Sub-Saharan Africa, the most malarious place on earth, presented such a daunting logistical challenge that the eradication campaign never really got under way there. And, even in countries that seemed highly amenable to spraying, problems arose. "The rich had houses that they didn't want to be sprayed, and they were giving bribes," says Socrates Litsios, who was a scientist with the W.H.O. for many years and is now a historian of the period. "The inspectors would try to double their spraying in the morning so they wouldn't have to carry around the heavy tanks all day, and as a result houses in the afternoon would get less coverage. And there were many instances of corruption with insecticides, because they were worth so much on the black market. People would apply diluted sprays even when they knew they were worthless." Typical of the logistical difficulties is what happened to the campaign in Malaysia. In Malaysian villages, the roofs of the houses were a thatch of palm fronds called atap. They were expensive to construct, and usually lasted five years. But within two years of DDT spraying the roofs started to fall down. As it happened, the atap is eaten by caterpillar larvae, which in turn are normally kept in check by parasitic wasps. But the DDT repelled the wasps, leaving the larvae free to devour the atap. "Then the Malaysians started to complain about bedbugs, and it turns out what normally happens is that ants like to eat bedbug larvae," McWilson Warren said. "But the ants were being killed by the DDT and the bedbugs weren't--they were pretty resistant to it. So now you had a bedbug problem." He went on, "The DDT spray teams would go into villages, and no one would be at home and the doors would be locked and you couldn't spray the house. And, understand, for that campaign to work almost every house had to be sprayed. You had to have eighty-per-cent coverage. I remember there was a malaria meeting in '62 in Saigon, and the Malaysians were saying that they could not eradicate malaria. It was not possible. And everyone was arguing with them, and they were saying, 'Look, it's not going to work.' And if Malaysia couldn't do it--and Malaysia was one of the most sophisticated places in the region--who could?"
At the same time, in certain areas DDT began to lose its potency. DDT kills by attacking a mosquito's nervous system, affecting the nerve cells so that they keep firing and the insect goes into a spasm, lurching, shuddering, and twitching before it dies. But in every population of mosquitoes there are a handful with a random genetic mutation that renders DDT nontoxic--that prevents it from binding to nerve endings. When mass spraying starts, those genetic outliers are too rare to matter. But, as time goes on, they are the only mosquitoes still breeding, and entire new generations of insects become resistant. In Greece, in the late nineteen-forties, for example, a malariologist noticed Anopheles sacharovi mosquitoes flying around a room that had been sprayed with DDT. In time, resistance began to emerge in areas where spraying was heaviest. To the malaria warriors, it was a shock. "Why should they have known?" Janet Hemingway, an expert in DDT resistance at the University of Wales in Cardiff, says. "It was the first synthetic insecticide. They just assumed that it would keep on working, and that the insects couldn't do much about it." Soper and the malariologist Paul Russell, who was his great ally, responded by pushing for an all-out war on malaria. We had to use DDT, they argued, or lose it. "If countries, due to lack of funds, have to proceed slowly, resistance is almost certain to appear and eradication will become economically impossible," Russell wrote in a 1956 report. "TIME IS OF THE ESSENCE because DDT resistance has appeared in six or seven years." But, with the administrative and logistical problems posed by the goal of eighty-per-cent coverage, that deadline proved impossible to meet.
In 1963, the money from Congress ran out. Countries that had been told they could wipe out malaria in four years--and had diverted much of their health budgets to that effort--grew disillusioned as the years dragged on and eradication never materialized. Soon, they put their money back into areas that seemed equally pressing, like maternal and child health. Spraying programs were scaled back. In those countries where the disease had not been completely eliminated, malaria rates began to inch upward. In 1969, the World Health Organization formally abandoned global eradication, and in the ensuing years it proved impossible to muster any great enthusiasm from donors to fund antimalaria efforts. The W.H.O. now recommends that countries treat the disease largely through the health-care system--through elimination of the parasite--but many anti-malarial drugs are no longer effective. In the past thirty years, there have been outbreaks in India, Sri Lanka, Brazil, and South Korea, among other places. "Our troubles with mosquitoes are getting worse," Spielman concludes in "Mosquito," "making more people sick and claiming more lives, millions of lives, every year."
For Soper, the unravelling of his dream was pure torture. In 1959, he toured Asia to check on the eradication campaigns of Thailand, the Philippines, Ceylon, and India, and came back appalled at what he had seen. Again and again, he found, countries were executing his strategy improperly. They weren't spraying for long enough. They didn't realize that unless malaria was ground into submission it would come roaring back. But what could he do? He had prevailed against gambiae in Brazil in the nineteen-thirties because he had been in charge; he had worked with the country's dictator to make it illegal to prevent an inspector from entering a house, and illegal to prevent the inspector from treating any open container of water. Jesse Hobbs tells of running into Soper one day in Trinidad, after driving all day in an open jeep through the tropical heat. Soper drove up in a car and asked Hobbs to get in; Hobbs demurred, gesturing at his sweaty shirt. "Son," Soper responded, "we used to go out in a day like this in Brazil and if we found a sector chief whose shirt was not wet we'd fire him." Killing mosquitoes, Soper always said, was not a matter of knowledge and academic understanding; it was a matter of administration and discipline. "He used to say that if you have a democracy you can't have eradication," Litsios says. "When Soper was looking for a job at Johns Hopkins--this would have been '46--he told a friend that 'they turned me down because they said I was a fascist.'" Johns Hopkins was right, of course: he was a fascist--a disease fascist--because he believed a malaria warrior had to be. But now roofs were falling down in Malaysia, and inspectors were taking bribes, and local health officials did not understand the basic principles of eradication--and his critics had the audacity to blame his ideas, rather than their own weakness.
It was in this same period that Rachel Carson published "Silent Spring," taking aim at the environmental consequences of DDT. "The world has heard much of the triumphant war against disease through the control of insect vectors of infection," she wrote, alluding to the efforts of men like Soper, "but it has heard little of the other side of the story--the defeats, the short-lived triumphs that now strongly support the alarming view that the insect enemy has been made actually stronger by our efforts." There had already been "warnings," she wrote, of the problems created by pesticides:
On Nissan Island in the South Pacific, for example, spraying had been carried on intensively during the Second World War, but was stopped when hostilities came to an end. Soon swarms of a malaria-carrying mosquito reinvaded the island. All of its predators had been killed off and there had not been time for new populations to become established. The way was therefore clear for a tremendous population explosion. Marshall Laird, who had described this incident, compares chemical control to a treadmill; once we have set foot on it we are unable to stop for fear of the consequences.
It is hard to read that passage and not feel the heat of Soper's indignation. He was familiar with "Silent Spring"--everyone in the malaria world was--and what was Carson saying?Of course the mosquitoes came back when DDT spraying stopped. The question was whether the mosquitoes were gone long enough to disrupt the cycle of malaria transmission. The whole point of eradication, to his mind, was that it got you off the treadmill: DDT was so effective that if you used it properly you could stop spraying and not fear the consequences. Hadn't that happened in places like Taiwan and Jamaica and Sardinia?
"Silent Spring" was concerned principally with the indiscriminate use of DDT for agricultural purposes; in the nineteen-fifties, it was being sprayed like water in the Western countryside, in an attempt to control pests like the gypsy moth and the spruce budworm. Not all of Carson's concerns about the health effects of DDT have stood the test of time--it has yet to be conclusively linked to human illness--but her larger point was justified: DDT was being used without concern for its environmental consequences. It must have galled Soper, however, to see how Carson effectively lumped the malaria warriors with those who used DDT for economic gain. Nowhere in "Silent Spring" did Carson acknowledge that the chemical she was excoriating as a menace had, in the two previous decades, been used by malariologists to save somewhere in the vicinity of ten million lives. Nor did she make it clear how judiciously the public-health community was using the chemical. By the late fifties, health experts weren't drenching fields and streams and poisoning groundwater and killing fish. They were leaving a microscopic film on the inside walls of houses; spraying every house in a country the size of Guyana, for example, requires no more DDT in a year than a large cotton farm does. Carson quoted a housewife from Hinsdale, Illinois, who wrote about the damage left by several years of DDT spraying against bark beetles: "The town is almost devoid of robins and starlings; chickadees have not been on my shelf for two years, and this year the cardinals are gone too; the nesting population in the neighborhood seems to consist of one dove pair and perhaps one catbird family. . . . 'Will they ever come back?' [the children]ask, and I do not have the answer." Carson then quoted a bird-lover from Alabama:"There was not a sound of the song of a bird. It was eerie, terrifying. What was man doing to our perfect and beautiful world?" But to Soper the world was neither perfect nor beautiful, and the question of what man could do to nature was less critical than what nature, unimpeded, could do to man. Here, from a well-thumbed page inserted in Soper's diaries, is a description of a town in Egypt during that country's gambiae invasion of 1943--a village in the grip of its own, very different, unnatural silence:
Most houses are without roofs. They are just a square of dirty earth. In those courtyards and behind the doors of these hovels were found whole families lying on the floor; some were just too weakened by illness to get up and others were lying doubled up shaking from head to foot with their teeth chattering and their violently trembling hands trying in vain to draw some dirty rags around them for warmth. They were in the middle of the malaria crisis. There was illness in every house. There was hardly a house which had not had its dead and those who were left were living skeletons, their old clothing in rags, their limbs swollen from undernourishment and too weak to go into the fields to work or even to get food.
It must have seemed to Soper that the ground had shifted beneath his feet--that the absolutes that governed his life, that countenanced even the most extreme of measures in the fight against disease, had suddenly and bewilderingly been set aside. "I was on several groups who evaluated malaria-eradication programs in some of the Central American countries and elsewhere," Geoffrey Jeffery recalls. "Several times we came back with the answer that with the present technology and effort it wasn't going to work. Well, that didn't suit Soper very much. He harangued us. We shouldn't be saying things like that!" Wilbur Downs, a physician who worked for the Rockefeller Foundation in Mexico in the fifties, used to tell of a meeting with Soper and officials of the Mexican government about the eradication of malaria in that country. Soper had come down from Washington, and amid excited talk of ending malaria forever Downs pointed out that there were serious obstacles to eradication--among them the hastened decomposition and absorption of DDT by the clays forming adobe walls. It was all too much for Soper. This was the kind of talk that was impeding eradication--the doubting, the equivocation, the incompetence, the elevation of songbirds over human life. In the middle of the meeting, Soper--ramrod straight, eyes afire--strode over to Downs, put both his hands around his neck, and began to shake.
Fred Soper ran up against the great moral of the late twentieth century--that even the best-intentioned efforts have perverse consequences, that benefits are inevitably offset by risks. This was the lesson of "Silent Spring," and it was the lesson, too, that malariologists would take from the experience with global eradication. DDT, Spielman argues, ought to be used as selectively as possible, to quell major outbreaks. "They should have had a strong rule against spraying the same villages again and again," he says. "But that went against their doctrine. They wanted eighty-per-cent coverage. They wanted eight out of ten houses year after year after year, and that's a sure formula for resistance." Soper and Russell once argued about whether, in addition to house spraying, malaria fighters should continue to drain swamps. Russell said yes; Soper said no, that it would be an unnecessary distraction. Russell was right: it made no sense to use only one weapon against malaria. Spielman points out that malaria transmission in sub-Saharan Africa is powerfully affected by the fact that so many people live in mud huts. The walls of that kind of house need to be constantly replastered, and to do that villagers dig mud holes around their huts. But a mud hole is a prime breeding spot for gambiae. If economic aid were directed at helping villagers build houses out of brick, Spielman argues, malaria could be dealt a blow. Similarly, the Princeton University malariologist Burton Singer says that since the forties it has been well known that mosquito larvae that hatch in rice fields--a major breeding site in southeast Asia--can be killed if the water level in the fields is intermittently drained, a practice that has the additional effect of raising rice yields. Are these perfect measures? No. But, under the right circumstances, they are sustainable. In a speech Soper presented on eradication, he quoted Louis Pasteur: "It is within the power of man to rid himself of every parasitic disease." The key phrase, for Soper, was "within the power." Soper believed that the responsibility of the public-health professional was to make an obligation out of what was possible. He never understood that concessions had to be made to what was practical. "This is the fundamental difference between those of us in public health who have an epidemiological perspective, and people, like Soper, with more of a medical approach," Spielman says. "We deal with populations over time, populations of individuals. They deal with individuals at a moment in time. Their best outcome is total elimination of the condition in the shortest possible period. Our first goal is to cause no outbreaks, no epidemics, to manage, to contain the infection." Bringing the absolutist attitudes of medicine to a malarious village, Spielman says, "is a good way to do a bad thing." The Fred Soper that we needed, in retrospect, was a man of more modest ambitions.
But, of course, Fred Soper with modest ambitions would not be Fred Soper; his epic achievements arose from his fanaticism, his absolutism, his commitment to saving as many lives as possible in the shortest period of time. For all the talk of his misplaced ambition, there are few people in history to whom so many owe their lives. The Global Malaria Eradication Programme helped eliminate the disease from the developed world, and from many parts of the developing world. In a number of cases where the disease returned, it came back at a lower level than it had been in the prewar years, and even in those places where eradication made little headway the campaign sometimes left in place a public infrastructure that had not existed before. The problem was that Soper had raised expectations too high. He had said that the only acceptable outcome for Global Eradication was global eradication, and when that did not happen he was judged--and, most important, he judged himself--a failure. But isn't the urgency Soper felt just what is lacking in the reasonableness of our contemporary attitude--in our caution and thoughtfulness and restraint? In the wake of the failure of eradication, it was popular to say that truly effective malaria control would have to await the development of a public-health infrastructure in poorer countries. Soper's response was, invariably: What about now? In a letter to a friend, he snapped, "The delay in handling malaria until it can be done by local health units is needlessly sacrificing the generation now living." There is something to admire in that attitude; it is hard to look at the devastation wrought by H.I.V. and malaria and countless other diseases in the Third World and not conclude that what we need, more than anything, is someone who will marshal the troops, send them house to house, monitor their every movement, direct their every success, and, should a day of indifference leave their shirts unsullied, send them packing. Toward the end of his life, Soper, who died in 1975, met with an old colleague, M. A. Farid, with whom he had fought gambiae in Egypt years before. "How do things go?" Soper began. "Bad!" Farid replied, for this was in the years when everyone had turned against Soper's vision. "Who will be our ally?" Soper asked. And Farid said simply, "Malaria," and Soper, he remembered, almost hugged him, because it was clear what Farid meant: Someday, when DDT is dead and buried, and the West wakes up to a world engulfed by malaria, we will think back on Fred Soper and wish we had another to take his place. |
Types of Memory
Memory has two major divisions:
Procedural memory concerns our memories of how to do things. Procedual memory guides the processes we perform and normally resides below the level of conscious awareness. When needed, procedural memories are automatically retrieved and used for both cognitive and motor skills (e.g. tying laces, driving a car, writing). Procedural memory is created through "procedural learning" (repeating a complex activity over and over again) and making all of the relevant neural systems work together automatically. Implicit (no concious awareness) procedural learning is essential to the development of any motor skill or cognitive activity.
Declarative Memory refers to memory which can be consicuously recalled such as facts and knowledge. Declarative memory is further divided into semantic memory and episodic memory. However, the two types of memory may interact. It may be the case that these are different types of information stored in the same system as opposed to being seperate systems themselves.
Semantic memory refers to the memory of meanings, understandings, and other concept-based knowledge unrelated to personal experiences. Basically, semantic memory is the knowledge of facts, such as knowing what 'x' is. Semantic memory is best described as "what we already understand" while episodic memory is best described as "what has happened to us". Semantic memory consists of networks of associations between concepts. Links between these concepts (nodes) are labelled in a semantic fashion, with there being at least two types of links in the Semantic Network, these being class membership and attributes.
As opposed to episodic memory, semantic memory is not affected by Amnesia which is a condition when one's memory is lost. However, semantic memory is affected by Agnosia (loss of knowledge- it can be the loss of ability to recognize objects, persons, sounds, shapes, or smells). Furthermore, it is unrelated to context and personal relevance.
Hierarchical Network Model (HNM)
The Hierarchical network model (HNM) by Collins & Quillian (1969), was the first systematic model of semantic memory. The model suggests that semantic memory is organised into a series of hierarchical networks, consisting of nodes and properties. A node is a major concept, such as 'animal, bird, canary'. A property, attribute or feature is, as expected, a property of that concept. For example 'has wings, is yellow'. The model is arranged as a hierarchy, with the more widely encompassing nodes stored on the higher levels. The underlying principle of the model is that of cognitive economy. Property information is stored as high up as possible to minimise the amount of information stored in semantic memory. This means that 'fish' is on a higher level than 'fresh water fish', which in turn is on a higher level than 'salmon'. We make inferential choices in semantic memory. For example the information that Picasso had knees is not stored in semantic memory. The knowledge that he was a human is and we know that humans have knees so we can infer that Picasso must have had knees.
When a concept is 'activated' in semantic memory, linked nodes are also 'activated' and relevant data is inferred. An example of this is people were faster to react to "a canary is yellow" than "a canary has wings". This illustrates that the closer together in the hierachy, the faster someone can identify concepts and their properties. The concept (canary) and the property (yellow) are stored at the same level, and are thus activated quickly, but canary and "can fly" are separated by one level, and so reaction time takes longer.
Collins implemented the hierarchical model of semantic memory in a computer programme that could understand basic text - the Teachable Language Comprehender (TLC).
Evidence in favour
Collins and Quillian (1969, 1972) found that when asking questions such as 'a canary is an animal' and 'a canary is a canary', the results were as predicted by the model in that the greater the semantic distance the slower reaction time.
Whilst the HNM had the right idea regarding semantic memory using a process of inferral, there are several issues with the model in general.
If falsification is taken as the means of measurement in reaction times, it becomes apparent that the greater the semantic distance between the concepts (nodes) the faster the reaction times. Collins and Quillian themselves found that 'a canary is a tulip' was faster to be rejected than 'a canary is a robin'. The model would suggest that for the tulip question the participant would have to search through the whole network before rejecting it and thus would be slower.
Conrad (1972) investigated whether hierarchical distance, or familiarity, was more influential in determining whether sentences were true or false. This was based on the notion that a slower reaction time for verification questions such as "a canary has skin", compared to "a canary sings", could be more down to the familiarity of the sentences. Conrad controlled the familiarity, and found that hierarchical distance between the subject and the property had little effect on verification time.
Another issue is that the model fails to explain typicality effects. The finding that people are faster to verify category inclusion for typical category members. For example, "a canary is a bird" is faster to verify than "a penguin is a bird", as a penguin is less typical/representative of the category of bird. This is supported by Rosch and Mervis (1975) who investigated the typicality ratings of fruits and found that oranges, apples, bananas and pears were rated as much more typical fruits than olives, tomatoes, coconuts and dates. Rips, Shoben and Smith (1973), found that verification times were faster for more typical or representative members, than for more atypical members of their category. This is called the typicality gradient.
Rosch (1973) showed that more typical members shared more characteristics associated with a category than atypical members such as 'a robin is a bird' registerd faster than 'a chicken is a bird'. This suggests that the concepts we use are much more loosely fitted to categories than the HNM proposes. Strong support for this came from McCloskey and Glucksberg (1978). They asked participants 30 'tricky' questions, such as "is a stroke a disease?". Answers were not unanimous among participants, and 11 months later many people had changed their answers to some of the questions. This shows how fuzzy memory can be.
A further issue is that the model only accounts for remembering sentences of a specific form (e.g a bird has wings)
Whilst the HNM is clearly flawed, it is worth remembering that it was the first systematic model of semantic memory, and it's influence over later models should not be dismissed.
Teachable Language Comprehender (TLC) Collins and Quillian (1969)
To comprehend text input by relating it to a pre-existing large semantic network (SN) representing rules already known about the world. Comprehension is identical with successful relation to input and learning is accomplished by incorporating any successfully comprehended rules into the SN.
Similar concepts are stored closer together than unrelated concepts
- Concepts are stores as local representations – each concept is stored as a single node and the nodes of related concepts are linked together in a hierarchal fashion.
- Links represent relationships between nodes – ‘is a’ or ‘has’.
- Cognitive economy – minimises the number of representations of a piece of information
- All or none
-Principle of cognitive economy is used in storage of properties for a concept. A property is stored at the highest possible node in the hierarchy so information can be deduced via inheritence for lower nodes e.g. 'has wings'.
Process of semantic analysis
- Intersection search: as soon as a word is parsed (broken down and analysed) it spreads like a 'plague', recording its original form and previous 'victim' so it is possible to retract. When it reaches a node it has already 'touched' this links the two nodes semantically (length of path= semantic distance).
- Semantic Interpretation of input corresponds to the set of linked words. E.g. for the phrase ‘the canary the shark bit had wings’, the semantic network is used to infer that the canary is the owner of the wings, and the shark is the one that bit, because ‘canary’ has ‘wings’ as a semantically linked property, and ‘shark’ has ‘biting’ as a semantically linked property.
- Syntax is only used to check the validity of interpretation. Any input that is not syntactically correct is rejected.
Classes of failure
- Because of its generalisation hierarchy structure, there was nowhere to put any abstract information that didn’t fit into the hierarchy.
- Often made false connections between subjects if they were too general and the syntax was too vague. E.g. ‘he hated the landlord so much that he moved into the house on Brunswick Street’ – TLC would incorrectly associate the landlord and the house.
- Although it is possible to comprehend episodic input using the semantic network, it doesn’t incorporate these episodes into the network itself, i.e. it can only learn semantic relationships
Studies on sentence varification times (Collins and Quillian 1969/72) show good support for the notion that reaction time increases as the semantic distance increases.
Problems for the TLC:
- Own data is inconsistent with model as RTs are faster the greater the semantic distance.
- Typicality effects; no associative strength attached to links shown by Rosch (1973). 'A robin is a bird' verified faster than 'a chicken is a bird' due to the fact that there is a difference in typicality between the two. Ratings of typicality were robin-bird (1.1) and chicken-bird (3.8) on a 1-7 rating scale.
- Alternative explanation for sentence verification= issue of typicality/representation. No evidence for cognitive economy (Conrad 1972).
It is unlikely that the precise representation chosen bears much resemblance to human semantic memory. However, TLC was hugely influential, first in demonstrating that it's possible to model SM and second, in influencing the development of consequent better models.
The Spreading Activation Model
Collins and Loftus (1975) developed the spreading activation model of semantic memory as a more complex answer to the HNM's criticisms. It suggests that concepts and nodes are linked together with different levels of conductivity. The more often the two concepts are linked, the greater conductivity. Thus the conductivity may be thought of as the criteria of the relation.
Collins and Loftus assumed that semantic memory is organised on the basis of semantic relatedness or semantic distance. In this model elements are linked in a conductive manner - the more two elements are activated together the greater their conductivity. This is modeled through shorter links between nodes. The shorter the link, the closer the semantic relation, and so the faster the brain will be at making the connection between the nodes. Furthermore, the longer a concept is accessed, the larger the spread of activation. When a concept is accessed activation spreads out from that node in all directions. The higher the conductivity the faster it spreads down that link. Whenever a person thinks hears or see a concept the appropriate node is activated. Like neurons, each intersection has a threshold and activation summates linearly from different inputs to the node.
The principle of weak cognitive economy is basically a revised version of Collins and Quillian's cognitive economy principle that allows information to be stored at a lower node in the hierarchy if the link has been explicit, even if already stored at a higher level. If relations are not stored explicitly it is still possible to infer them using hierarchical information.
Collins and Loftus claim that there are different types of links including: -Class membership (a cat is a mammal), -Subordinate (a cat has fur), -Prediction (game-play-people) -Exclusion (a whale is not a fish). Collins and Loftus suggest that connections made are not necessarily logical, rather based on personal experience.
The model can explain the familiarity effect, the typicality effect, and direct concept-property associations.
S.A.M is supported by studies of priming in which there is an improvement in speed or accuracy to respond to a stimulus when it is seen to proceed a semantically related concept. Mackay (1973)demonstrated for example how prior context can remove any disambiguation from a phrase (e.g. he walked towards the bank) because of these interconnected units of information. Meyer and Schvaneveldt (1971) found that when words are related, reaction times are quicker. They asked participants if both words in a pair were words or non words. Participants answered "yes" much quicker when the words were related (e.g. bread and butter) compared to when they were not related (e.g. bread and coat). If the words are related activation from the first word is spread to the second word making the association much faster than if they are not related.
However the disadvantage is that the theory is unable to predict much as it is based on the individual. It handles everything and makes very few predictions which are open to empirical testing making it very difficult to falsify.. The model also fails to consider how episodic knowledge or non-propositional knowledge could be stored. There are so many possible parameters to the system that it is possible to fit almost any empirical data anyway. Despite its neurological plausibility it is not sufficiently constrained enough to allow it to be implemented reliably.
SAM however, offers many strengths as a model, as it has clear face validity. It seems a plausible model and gives a good foundation of how the semantic network is built up initially.
For information to be used in a task like recognition, it must first be activated and then inspected. When information is in the LTM, but not currently in the WM, activation must spread to it.
The Fan Effect
The Fan effect causes interference in semantic memory. The more facts that are associated with a 'node', the slower the activation spreads from it, as a node has a fixed capacity for emitting activation. Therefore if there are more links to that node then more time must be taken in order to activate all the links, although this can be sped up if the links are often used and so there is more immediate association.
Anderson (1974) asked participants to learn sentences comprising of a subject and location with a relation between them. For example:
1. The Doctor is in the bank
2. The fireman is in the park
3. The lawyer is in the Church
4. The lawyer is in the park
Participants were then given a speed recognition task. They were asked to indicate when they recognised a learnt sentence (the target) amongst other sentences of similar nature (the distractors) An example of a distractor may be "the Doctor is in the park". Anderson found participants reaction time was faster when there were less shared facts. Reaction time for unique sentences -e.g. "the Doctor is in the bank" - was 1.11 seconds compared to when the location and person appeared in two sentences -e.g. "the lawyer is in the park" - which was at 1.22 seconds. Thus the more facts are associated with a node the slower the activation spreads from it - this is the fan effect. This has implications- the more you know the slower you get?! This may be the case but we can speed this up consciously using procedural memory.
Semantic nodes are very subjective as people's schemas differ significantly. For example, the word 'apple' might eliit the colour 'green' to one person, yet 'red' to another. This could slow down the Fan effect due to the collaboration of knowledge that others may chose differently; hence slower the activiation.
Adaptive Control of Thought (ACT*)Declarative Memory
ACT* was built upon the TLC and SAM models of semantic memory. It maintained the idea of “semantic networks” but suggested it was “activation” that was key to semantic knowledge and memory. The ACT* model was the first complete model of human cognition, (a challenging task). It therefore has a highly complex architecture which allows it to learn. This type of knowledge is called declarative; the knowledge of facts and information. ACT* also suggests that human knowledge can be procedural; knowledge we hold in order to perform automatic actions such as driving.
The ACT* therefore suggests that memories must be activated from source nodes in the Working Memory. Studies have highlighted that activation takes place automatically, that is, it requires no conscious awareness. The ACT* Throy of fact recognition proposed that items in the LTM remain permanently but cannot be accessed directly unless they are 'activated'. This model viewed Working Memory differently to how it had previously been viewed (Baddeley and Hitch) whereby 'source nodes' could be located anywhere throughout the brain and all those that are activated at any one time make up Working Memory. Being a complex system the ACT* is most easily represented by the “lightbulb analogy”. If you think of a floor of interconnected lightbulbs, most of which are off, some are dim (partially activated) and some will be lit brightly (fully activated). At different times, different sections of light bulbs will be turned off and on. This is supposed to represent the idea that activation is a continuous function rather than an all-or-none action.
Conclusions on Semantic Memory
There is no doubt that we have some sort of SN in our brains developing through experience. Semantic memory plays a crucial role in almost any cognitive activity and it is very likely that some sort of speading activation is involved in accessing this system. However, the system is not constrained enough to allow us to decide which complex model is closest to the truth. We don't have a good understandingof the way other parts of the system work. Semantic memories are learned through experiences so are idiosyncratic, that is, particular to the individual. Finally, Cognitive neuroscience may provide further insight into regional brain activity during SM tasks.
Schemas are a set of rules on how to behave which can be applied to a situation. We use schemas every day, for instance when we go to a lecture. Our lecture schema tells us that we have to find out where our lecture is, how to get there, get there, enter the lecture theatre, find a seat and so on. Bartlett introduced the concept. Schemas are automatic and unconscious processes that occur on a daily basis. Our schemata tend to be long lasting and are not easily changed. If we are introduced to information that contradicts our schemata we tend to assume that the new information is unique or different rather than believing our schemata are faulty. This links to Piaget's concepts of accomodation and assimilation- accomodation is the process of adapting existing schemata to fit with new information, whereas assimilation involves modifying new information to fit with our preexisting schemas.
Language, Scripts and Frames
The given-new contract(Clark 1977) describes how, in a conversation, a speaker should provide "not too much or too little" in terms of establishing the context.
1. The information given by the speaker should provide the appropriate context from which new information can be given. To illustrate this, in the example "Jamie went to the shops" - it is "given" that we are talking about "Jamie", as this has been established at the forefront of the utterance. This establishment of context provides a platform for "new" information concerning the fact of where Jamie went - "to the shops".
2.In establishing this context, the speaker should not give more information than needed- only the amount necessary to establish the context of conversation. This forms the basis of cost-effective communication. An example of this is being asked directions by someone foreign, you accomodate thier lack of knowledge and give simple directions in detail. In contrast, when giving directions to someone local you may assume they know certain landmarks/places so will give appropriate level of information.
3.Enough appropriate information must be given in order for the listener to be able to make "bridging inferences", this is where the listener is able to refer back to previous elements of the conversation in order to infer what is being discussed at later parts of the conversation. To illustrate this, if we extend the previous example to: "Jamie went to the shops with Luke, he bought a loaf a bread", the listener can infer that "he" refers to Jamie - by recognizing that Jamie is the main subject/actor in the utterance.
4.Lastly, a speaker may include more or less specific information depending upon the knowledge base of the listener. For example, if I bring up the subject of the components of working memory to someone who takes Art, I will most likely provide a generic description of the matter in hand - as I am aware that my listener will only be able to make bridging inferences based on their general knowledge.
The necessity of the given-new contract was highlighted in a study by Bransford and Johnson(1972). They asked participants to read this paragraph:
'The procedure is actually quite simple. First you arrange things into different groups. Of course, one pile may be sufficient depending on how much there is to do. If you have to go somewhere else due to lack of facilities, that is the next step, otherwise you are pretty well set. It is important not to overdo things. That is, it is better to do a few things at once rather than too many. In the short run this may not seem too important, but complications can easily arise. A mistake can be expensive as well.'
When participants were told beforehand that the paragraph was about laundry, they were able to understand it and recall what was said. If participants were not told anything about the passage, they found it difficult to understand. This suggests that having the memory schemata activated led to comprehension and retention of the information.
Frames (Minsky 1976)
Frames consist of schema for organising information about a single concept. It has 'slots' that can be filled with 'variables' which can be compulsory or optional. If no information is given explicitly, optional variables may be filled with 'default' values. If all of these slots are filled, the frame is instantiated. Problems: 'bare' frames don't allow the use of context so needs a more global structure that is able to account for meaningful 'chunks' of life in which frames can be placed. This corresponds to Barlett's idea of schemas but they are recently described as 'scripts'.
Frames don't need to store extra information. Bower, Black and Turner (1979) asked subject to list actions normally taken in certain situations such as going to the doctors. They found lots of commonality between scripts and they also found that the more scripts a participant read, the larger the intrusion of unstated generic script material in recall.
Shank and Abelson (1970-82) Natural Language Understanding
SAM - Script Applier Mechanism
A computer program developed to create a script of default knowledge which is normally taken for granted in a conversation. The computer creates these scripts so that it can 'understand' the conversation as it requires the scripts human beings would make normally. When given a statement, SAM paraphrases it and can then answer questions on the original statement. The script provides the default information that allows the the computer to make necessary inferences to understand the statement. It is not possible to create a script for every eventuality, so sometimes we make inferences by understanding what is a likely outcome/reaction.
PAM - Plan Applier Mechanism
PAM was created by Schank and Abelson (1970-82). Whenever an inference needs to be made, PAM tried to find a link using a repertoire of possible plans. It is not perfect, as it generates ENDLESS possibilities, but it is clear we use our knowledge of what WE might do, to infer what OTHERS might do.
It has been argued that SAM is a much stonger concept as it is clear that we constantly face situations where not only our own knowledge is relevant but also the knowledge of others; the ability to put ourselves in 'someone elses shoes'.
Does knowledge help?
De Groot (1965) showed subjects a possition on chessboard before sweeping it off and asking for it to be replicated. 10 year old experts did better that 20 year old beginners and there was no difference when the possitions were scrambled suggesting knowledge was instrumental in reconstruction. Similarly, Voss et al 1978 found participants with high baseball knowledge did significantly better in recall test after reading a baseball passage than compared to low knowledge group.
Summary on expertise and links to language and memory
Experts use knowledge to develop abstract, highly specialised mechanisms for systematically encoding and retrieving meaningful patterns from LTM. It allows experts to anticipate information needed for familiar tasks and stores new information in a format that facilitates retrieval.
We are expert in our own social world- we have developed mechanisms for systematically encoding meaningful language information. Spreading Activation uses existing declarative memory structures to automatically disambiguate utterances, providing a mechanism in which schemas/scripts may be instantiated in the human brain. |
Extrahepatic manifestations of hepatitis C.It is well established that chronic hepatitis Chronic hepatitis
Long lasting inflammation of the liver due to viruses or other causes.
Mentioned in: Tube Compression of the Esophagus and Stomach
chronic hepatitis C virus (HCV HCV
hepatitis C virus
HCV 1 Hepatitis C virus, see there 2. Human coronavirus. See Coronavirus. ) infection may lead to the development of progressive liver disease Liver Disease Definition
Liver disease is a general term for any damage that reduces the functioning of the liver.
The liver is a large, solid organ located in the upper right-hand side of the abdomen. and cirrhosis of liver. For example, chronic HCV leads to cirrhosis in 20% of patients over 20 to 30 years. (1) This has lead to considerable focus on the hepatic manifestations of the HCV. However, there are research studies indicating that HCV is associated with various extrahepatic ex·tra·he·pat·ic
Originating or occurring outside the liver. manifestations including mixed cryoglobulinemia mixed cryoglobulinemia Mixed polyclonal-polyclonal cryoglobulinemia Nephrology A form of cryoglobulinemia characterized by a IgG & IgM, ± IgA cryoglobulins, evoked by to rheumatoid arthritis, SLE, Sjögren syndrome, EBV, CMV, subacute bacterial (MC), membranoproliferative glomerulonephritis mem·bra·no·pro·lif·er·a·tive glomerulonephritis
Chronic glomerulonephritis characterized by mesangial cell proliferation, increased lobular separation of glomeruli, thickening of glomerular capillary walls, and low serum levels of complement. (MPGN MPGN Membranoproliferative glomerulonephritis, see there ), non-Hodgkin lymphoma, Sjogren syndrome Sjögren syndrome Rheumatology An autoimmune disorder more common in older ♀, associated with rheumatoid arthritis, SLE, scleroderma, polymyositis, and other connective tissue diseases Clinical Dry eyes, dry mouth Lab SS-A, SS-B antibodies Management , porphyria Porphyria
comes in a winter storm to show her devotion, and her lover strangles her with her own tresses. [Br. Poetry: Browning Porphyria’s Lover in Magill IV, 247]
See : Love, Unrequited cutaneous cutaneous /cu·ta·ne·ous/ (ku-ta´ne-us) pertaining to the skin.
Of, relating to, or affecting the skin.
Pertaining to the skin. tarda, lichen planus Lichen Planus Definition
Lichen planus is a skin condition of unknown origin that produces small, shiny, flat-topped, itchy pink or purple raised spots on the wrists, forearms or lower legs, especially in middle-aged patients. , leukocytoclastic vasculitis, and various endocrine and neurologic manifestations. (1-5) Approximately 40% of patients with HCV have cryoglobulinemia with or without specific extrahepatic manifestation. (4-6) It has been suggested that focus on comorbid extrahepatic disorders may lead to better detection of HCV, improvement in these disorders due to antiviral treatment of HCV, and early diagnosis and treatment of these disorders. (7) More importantly, these manifestations can complicate the course and treatment of HCV-related liver diseases. For example, MC associated with HCV has been linked to increased fibrosis and cirrhosis of liver. Furthermore, it is possible that drug treatment of extrahepatic disorders may cause hepatotoxicity hepatotoxicity (hepˑ··tō·t including liver failure. Hence, recognition of type and magnitude of extrahepatic disorders, including their pathophysiology pathophysiology /patho·phys·i·ol·o·gy/ (-fiz?e-ol´ah-je) the physiology of disordered function.
1. , are very important to aid in awareness, early diagnosis, and specifically targeted therapy for HCV-related extrahepatic diseases.
MC is an autoimmune disease associated with chronic HCV. MC has frequently been implicated im·pli·cate
tr.v. im·pli·cat·ed, im·pli·cat·ing, im·pli·cates
1. To involve or connect intimately or incriminatingly: evidence that implicates others in the plot.
2. in increased fibrosis or cirrhosis in chronic HCV-infected patients. (6,8-11) The pathophysiologic mechanism of MC can also explain other extrahepatic syndromes such as MPGN, non-Hodgkin lymphoma, and leukocytoclastic vasculitis. (12-14) It has been postulated that HCV infects circulating B lymphocytes and can stimulate them to produce monoclonal immunoglobulin M (IgM) rheumatoid factor (Type II) and polyclonal polyclonal /poly·clo·nal/ (-klon´'l)
1. derived from different cells.
2. pertaining to several clones.
derived from different cells; pertaining to several clones. IgM rheumatoid factor (Type III) leading to mixed cryoglobulinemia. (15-16) These IgM rheumatoid factors bind with anti-HCV immunoglobulin G (IgG) or to the IgG-HCV immune complex leading to a deposition of circulating immune complexes in small- or medium-sized blood vessels ultimately resulting in essential MC, MPGN, MPGN and leukocytoclastic vasculitis. (17-18)
The exact nature of the relationship between MC and progressive liver disease remains unclear. It has been hypothesized that the association between cirrhosis and MC may be due to the association of cryoglobulin cryoglobulin /cryo·glob·u·lin/ (-glob´u-lin) an abnormal globulin that precipitates at low temperatures and redissolves at 37° C.
n. with longstanding HCV infection and older age. (10) However, fibrosis progression rate is clearly shown to be related to duration of infection rather than MC. (19) Furthermore, Siagris et al reported no significant correlation of duration of infection or older age with prevalence of MC in HCV-infected patients. (10) Hence, it is still not clear from the findings that the development of MC is cause or effect of progressive liver fibrosis or cirrhosis in chronic HCV infection. Further prospective longitudinal studies with histologic evaluation of liver biopsies comparing chronic HCV patients with or without cryoglobulinemia are warranted to clarify molecular and clinical relationship between MC and liver fibrosis.
Various neurologic manifestations have been reported in patients with chronic hepatitis C. (20-21) The most common peripheral neurologic deficit is mononeuropathy multiplex, which is clearly related to cryoglobulinemia, vasculitis Vasculitis Definition
Vasculitis refers to a varied group of disorders which all share a common underlying problem of inflammation of a blood vessel or blood vessels. The inflammation may affect any size blood vessel, anywhere in the body. , and thrombosis. (20-21) Other neurologic disorders associated with HCV include acute inflammatory syndromes, such as encephalitis encephalitis (ĕnsĕf'əlī`təs), general term used to describe a diffuse inflammation of the brain and spinal cord, usually of viral origin, often transmitted by mosquitoes, in contrast to a bacterial infection of the meninges , and encephalomyelitis encephalomyelitis /en·ceph·a·lo·my·eli·tis/ (en-sef?ah-lo-mi?e-li´tis) inflammation of the brain and spinal cord.
acute disseminated encephalomyelitis . It has been postulated that, since HCV belongs to a family of Flaviviridae, which is characterized by neurotropism neurotropism /neu·rot·ro·pism/ (ndbobr-rot´ro-pizm)
1. the quality of having a special affinity for nervous tissue.
2. , it has the ability to invade the nervous system. HCV RNA RNA: see nucleic acid.
in full ribonucleic acid
One of the two main types of nucleic acid (the other being DNA), which functions in cellular protein synthesis in all living cells and replaces DNA as the carrier of genetic is seen in affected brain tissue or cerebrospinal fluid (22-24) of patients with certain acute inflammatory central nervous system syndromes. Negative-strand HCV was identified in brain tissue of two patients with post-transplantation recurrent HCV and three of six autopsied patients with chronic HCV infection. (25-26) These findings support the idea that HCV may cause neurologic deficits or symptoms directly by invading and replicating in the nervous system. However, due to lack of suitable methodologies for in situ detection, it is unclear whether HCV has the same abilities of infection of the nervous system as its class. As of yet, the pathogenesis of nervous system involvement remains unclear because it can be attributed to in situ HCV replication, toxicity of HCV-encoded proteins, or immune-mediated antiviral response. Further basic virologic and clinical studies will be necessary to fully understand the viral interaction with the nervous system and the resulting neurologic symptoms.
Type II diabetes Type II diabetes
Type II diabetes is the most common form of diabetes and usually appears in middle aged adults. It is often associated with obesity and may be delayed or controlled with diet and exercise.
Mentioned in: Diabetic Ketoacidosis mellitus (DM) has been implicated as an important endocrine manifestation of chronic HCV. (27-30) Even though there is no formal proof for the infectious disease model for type II DM, it has been hypothesized that HCV mediates type II DM in genetically susceptible individuals. Thuluvath et al provided support for both racial and environmental factors for infectious disease model for the development of DM in chronic HCV. (31) They reported an increased prevalence of type II DM in black patients with HCV compared with controls matched for race, severity of liver disease, and body mass index. This finding is significant since both HCV infection and black race have been reported as independent risk factors for type II DM. (32-33) It has been noted that HCV-positive diabetics have beta-islet cell dysfunction with decreased C-peptide levels and limited acute insulin responses. (11,32) In addition, an improvement of glucose metabolism have been observed after antiviral therapy in patients with chronic HCV even though there is no animal or in vitro models to test the hypothesis that HCV directly damage beta islet cells or disturbs their synthetic function. (11) Therefore, further genetic, animal model, tissue culture media, and mechanistic studies are needed to determine the role of HCV on glucose metabolism and pancreatic disease that leads to type II DM. In addition, prospective longitudinal studies are needed in chronic HCV patients on antiviral therapy to determine whether DM is a reversible process in patients who clear HCV infection.
In conclusion, HCV may lead to various extrahepatic manifestations. The majority of extrahepatic manifestations can be explained in terms of an immune-mediated syndrome related to MC even though the exact mechanism by which HCV may result in MC is still not clear. Furthermore, MC is associated with increased fibrosis and cirrhosis. Hence, early treatment of HCV infection in patients with MC may decrease the rate of fibrosis and may prevent cirrhosis including its complications. Mechanisms of other extrahepatic manifestations such as neurologic manifestations and DM related to chronic HCV are unclear and necessitate further research to clarify pathophysiology and clinical course of extrahepatic manifestations of chronic HCV.
1. Liang TJ, Rehermann B, Seeff LB, et al. Pathogenesis, natural history, treatment and prevention of hepatitis C. AnnIntern Med 2000;132:296-305.
2. Sharara AI, Hunt CM, Hamilton JD. Hepatitis C Ann Intern Med 1996;125:658-668.
3. Gumber SC, Chopra SC, Hepatitis C. A multifaceted disease, review of extrahepatic manifestations. Ann Intern Med 1995;123:615-620.
4. Pawlotsky JM, Ben Yahia M, Andre C, et al. Immunological disorders in C virus chronic active hepatitis chronic active hepatitis 1. Obsolete term. See Chronic hepatitis2. Chronic viral hepatitis : a prospective case-control study. Hepatology 1994;19:841-848.
5. Clifford BD, Donahue D, Smith L, et al. High prevalence of serological serological
pertaining to or emanating from serology.
one involving examination of blood serum usually for antibody. markers of autoimmunity in patients with chronic hepatitis C. Hepatology 1995;21:613-619.
6. Lunel F, Musset L, Cacoub P, et al. Cryoglobulinemia in chronic liver disease Chronic liver disease is a liver disease of slow process and persisting over a long period of time, resulting in a progressive destruction of the liver.
It includes amongst others:
7. El-Serag HB, Hampel H, Yeh C, et al. Extrahepatic manifestations of hepatitis C among United State male veterans. Hepatology 2002;36:1439-1445.
8. Schmidt WN, Stapleton JT, La Brecque DR, et al. Hepatitis C virus (HCV) infection and cryoglobulinemia: analysis of whole blood and plasma HCV-RNA concentrations and correlation with liver histology. Hepatology 2000;31:737-743.
9. Kayali Z, Buckwold VE, Zimmerman B, et al. Hepatitis C, Cryoglobulinemia, and cirrhosis: A meta-analysis. Hepatology 2002;36:978-985.
10. Siagris D, Christofidou M, Tsamandas A, et al. Cryoglobulinemia and progression of fibrosis in chronic HCV infection: cause or effect? Journal of Infect 2004;49:236-241.
11. Akriviadis EA, Xanthakis I, Navrozidou CH, et al. Prevalence of cryoglobulinemia in chronic hepatitis C virus infection and response to treatment with interferon-a. J Clin Gastroenterol 1997;25:612-620.
12. Pascual M, Perrin L, Giostra E, et al. Hepatitis C virus in patients with cryoglobulinemia type II. J Infect Dis 1990;162:569-570.
13. Dammacco F, Sansonno D. Antibodies to hepatitis C virus in essential mixed cryoglobulinemia. Clin Exp Rheum rheum (rldbomacm) any watery or catarrhal discharge.
A watery or thin mucous discharge from the eyes or nose.
any watery or catarrhal discharge. 1991;9:621-624.
14. Agnello V, Chung RT, Kaplan LM. A role for hepatitis C virus infection in type II cryoglobulinemia. N Engl J Med 1992;327:1490-1495.
15. Muller HM, Pfaff E, Goeser T, et al Peripheral blood leukocytes serve as a possible extrahepatic site for hepatitis C virus replication. J Gen Virol 1993;74:669-676.
16. Mieschar PA, Huang YP, Izui S. Type II cryoglobulinemia. Semin Hematol 1995;32:80-85.
17. Willson RA. Extrahepatic manifestations of chronic viral hepatitis. Am J Gastroenterology 1997;92:4-17.
18. Agnello V. The etiology and pathophysiology of mixed cryoglobulinemia secondary to hepatitis C virus infection. Springer Semin Immunopathol 1997;19:111-129.
19. Poynard T, Ratziu V, Benmanov Y, et al. Fibrosis in patients with chronic hepatitis C: detection and significance. Semin Liver Dis 2000;20:47-55.
20. Heckmann J, Kayser C, Heuss D, et al. Neurological manifestations of chronic hepatitis C. J Neurol 1999;246:486-491.
21. Tembl J, Ferrer J, Sevilla M, et al. Neurologic complications associated with Hepatitis C virus infection. Neurology 1999;53:861-864.
22. Bolay H, Soylemezoglu F, Nurlu G, et al. PCR PCR polymerase chain reaction.
polymerase chain reaction
Polymerase chain reaction (PCR) detected hepatitis C virus genome in the brain of a case with progressive encephalomyelitis with rigidity. Clin Neurol Neurosurg 1996;98:305-308.
23. Fujita H, Chuganji Y, Yaginuma M, et al. Acute encephalitis and hepatitis C virus infection: case report: acute encephalitis immediately prior to acute onset of hepatitis C virus infection. J. Gastroenterology Hepatol 1999;14:1129-1131.
24. Sacconi S, Salviati L, Merelli E. Acute disseminated encephalomyelitis acute disseminated encephalomyelitis
A diffuse inflammation of the brain and spinal cord usually caused by a perivascular hypersensitivity response. associated with hepatitis C virus infection. Arch Neurol 2001;58:1679-1681.
25. Vargus H, Laskus T, Radkowski M, et al. Detection of hepatitis C virus sequences in brain tissue obtained in recurrent hepatitis C after liver transplantation. Liver Transplantation 2002;8:1014-1019.
26. Radkowski M, Wilkinson J, Nowicki M, et al. Search for hepatitis C virus negative-strand RNA sequences and analysis of viral sequences in the central nervous system: Evidence of replication. J. Virol 2002;76:600-608.
27. Fraser GM, Harman I, Meller N, et al. Diabetes Mellitus is associated with chronic hepatitis C but not chronic hepatitis B virus infection. Isr J Med Sci 1996;32:526-530.
28. Knobler H, Schihmanter R, Zifroni A, et al. Diabetes Mellitus is associated with chronic hepatitis C virus infection. Mayo Clin Proc 2000;75:355-359.
29. Zein zein
the principal protein in maize. Has low nutritive value, being deficient in lysine and tryptophan. NN, Abdulkarim AS, Wiesner RH, et al. Prevalence of diabetes mellitus in patients with end-stage liver cirrhosis due to hepatitis C, alcohol or cholestatic disease. J Hepatol 2000;32:209-17.
30. Mehta SH, Brancati FL, Sulkowski MS, et al. Prevalence of type 2 diabetes mellitus Type 2 diabetes mellitus
One of the two major types of diabetes mellitus, characterized by late age of onset (30 years or older), insulin resistance, high levels of blood sugar, and little or no need for supple-mental insulin. among persons with hepatitis C virus infection in the United States. Ann Intern Med 2000;133:592-599.
31. Thuluvath PJ, John PR. Association between hepatitis C, diabetes mellitus, and race: A case-control study. Am J Gastroenterol 2003;98:438-441.
32. Caronia S, Taylor K, Pagliaro L, et al. Further evidence for an association between non-insulin-dpendent diabetes mellitus and chronic hepatitis C virus infection. Hepatology 1999;30:1059-1063.
33. Mason A, Lau J, Hoang N, et al. Association of diabetes mellitus and chronic hepatitis C virus infection. Hepatology 1998;28:328-333.
The employees of the Southern Medical Journal wish to extend our heartfelt sympathies to the victims of Hurricane Katrina and Hurricane Rita.
We would also like to thank the members of the medical profession who assisted in the relief efforts throughout the Gulf Coast. Your professionalism, courage and tenacity have been an inspiration. Your compassion has reminded us of the true meaning of the avocation to which we have devoted our lives.
The purpose of human life is to show compassion and the will to help others. --Albert Schweitzer
Smruti R. Mohanty, MD
From the Center for Liver Diseases, Department of Medicine, Section of Gastroenterology, the University of Chicago, Chicago, IL.
Reprint Requests to Smruti R. Mohanty, MD, Assistant Professor of Medicine, Center for Liver Diseases, Department of Medicine, Section of Gastroenterology, the University of Chicago, 5841 South Maryland Avenue, MC 7120, Chicago, IL 60637. Email: email@example.com
Accepted April 19, 2005 |
Henry Gray (18211865). Anatomy of the Human Body. 1918.
4c. The Fore-brain or Prosencephalon
The fore-brain or prosencephalon consists of: (1) the diencephalon, corresponding in a large measure to the third ventricle and the structures which bound it; and (2) the telencephalon, comprising the largest part of the brain, viz., the cerebral hemispheres; these hemispheres are intimately connected with each other across the middle line, and each contains a large cavity, named the lateral ventricle. The lateral ventricles communicate through the interventricular foramen with the third ventricle, but are separated from each other by a medial septum, the septum pellucidum; this contains a slit-like cavity, which does not communicate with the ventricles.
FIG. 715 Mesal aspect of a brain sectioned in the median sagittal plane. (See enlarged image)
The Diencephalon.The diencephalon is connected above and in front with the cerebral hemispheres; behind with the mid-brain. Its upper surface is concealed by the corpus callosum, and is covered by a fold of pia mater, named the tela chorioidea of the third ventricle; inferiorly it reaches to the base of the brain.
The diencephalon comprises: (1) the thalamencephalon; (2) the pars mamillaris hypothalami; and (3) the posterior part of the third ventricle. For descriptive purposes, however, it is more convenient to consider the whole of the third ventricle and its boundaries together; this necessitates the inclusion, under this heading, of the pars optica hypothalami and the corresponding part of the third ventriclestructures which properly belong to the telencephalon.
The Thalamencephalon.The thalamencephalon comprises: (1) the thalamus; (2) the metathalamus or corpora geniculata; and (3) the epithalamus, consisting of the trigonum habenulæ, the pineal body, and the posterior commissure.
The Thalami (optic thalamus) (Figs. 716,717) are two large ovoid masses, situated one on either side of the third ventricle and reaching for some distance behind that cavity. Each measures about 4 cm. in length, and presents two extremities, an anterior and a posterior, and four surfaces, superior, inferior, medial, and lateral.
The posterior extremity is expanded, directed backward and lateralward, and overlaps the superior colliculus. Medially it presents an angular prominence, the pulvinar, which is continued laterally into an oval swelling, the lateral geniculate body, while beneath the pulvinar, but separated from it by the superior brachium, is a second oval swelling, the medial geniculate body.
The superior surface is free, slightly convex, and covered by a layer of white substance, termed the stratum zonale. It is separated laterally from the caudate nucleus by a white band, the stria terminalis, and by the terminal vein. It is divided into a medial and a lateral portion by an oblique shallow furrow which runs from behind forward and medialward and corresponds with the lateral margin of the fornix; the lateral part forms a portion of the floor of the lateral ventricle, and is covered by the epithelial lining of this cavity; the medial part is covered by the tela chorioidea of the third ventricle, and is destitute of an epithelial covering. In front, the superior is separated from the medial surface by a salient margin, the tænia thalami, along which the epithelial lining of the third ventricle is reflected on to the under surface of the tela chorioidea. Behind, it is limited medially by a groove, the sulcus habenulæ, which intervenes between it and a small triangular area, termed the trigonum habenulæ.
The medial surface constitutes the upper part of the lateral wall of the third ventricle, and is connected to the corresponding surface of the opposite thalamus by a flattened gray band, the massa intermedia (middle or gray commissure). This mass averages about 1 cm. in its antero-posterior diameter: it sometimes consists of two parts and occasionally is absent. It contains nerve cells and nerve fibers; a few of the latter may cross the middle line, but most of them pass toward the middle line and then curve lateralward on the same side.
The lateral surface is in contact with a thick band of white substance which forms the occipital part of the internal capsule and separates the thalamus from the lentiform nucleus of the corpus striatum.
Structure.The thalamus consists chiefly of gray substance, but its upper surface is covered by a layer of white substance, named the stratum zonale, and its lateral surface by a similar layer termed the lateral medullary lamina. Its gray substance is incompletely subdivided into three partsanterior, medial, and lateralby a white layer, the medial medullary lamina. The anterior part comprises the anterior tubercle, the medial part lies next the lateral wall of the third ventricle while the lateral and largest part is interposed between the medullary laminæ and includes the pulvinar. The lateral part is traversed by numerous fibers which radiate from the thalamus into the internal capsule, and pass through the latter to the cerebral cortex. These three parts are built up of numerous nuclei, the connections of many of which are imperfectly known.
FIG. 718 Coronal section of brain through intermediate mass of third ventricle. (See enlarged image)
Connections.The thalamus may be regarded as a large ganglionic mass in which the ascending tracts of the tegmentum and a considerable proportion of the fibers of the optic tract end, and from the cells of which numerous fibers (thalamocortical) take origin, and radiate to almost every part of the cerebral cortex. The lemniscus, together with the other longitudinal strands of the tegmentum, enters its ventral part: the thalamomammillary fasciculus (bundle of Vicq dAzyr), from the corpus mammillare, enters in its anterior tubercle, while many of the fibers of the optic tract terminate in its posterior end. The thalamus also receives numerous fibers (corticothalamic) from the cells of the cerebral cortex. The fibers that arise from the cells of the thalamus form four principal groups or stalks: (a) those of the anterior stalk pass through the frontal part of the internal capsule to the frontal lobe; (b) the fibers of the posterior stalk (optic radiations) arise in the pulvinar and are conveyed through the occipital part of the internal capsule to the occipital lobe; (c) the fibers of the inferior stalk leave the under and medial surfaces of the thalamus, and pass beneath the lentiform nucleus to the temporal lobe and insula; (d) those of the parietal stalk pass from the lateral nucleus of the thalamus to the parietal lobe. Fibers also extend from the thalamus into the corpus striatumthose destined for the caudate nucleus leave the lateral surface, and those for the lentiform nucleus, the inferior surface of the thalamus.
The medial geniculate body (corpus geniculatum mediale; internal geniculate body; postgeniculatum) lies under cover of the pulvinar of the thalamus and on the lateral aspect of the corpora quadrigemina. Oval in shape, with its long axis directed forward and lateralward, it is lighter in color and smaller in size than the lateral. The inferior brachium from the inferior colliculus disappears under cover of it while from its lateral extremity a strand of fibers passes to join the optic tract. Entering it are many acoustic fibers from the lateral lemniscus. The medial geniculate bodies are connected with one another by the commissure of Gudden, which passes through the posterior part of the optic chiasma.
The lateral geniculate body (corpus geniculatum laterale; external geniculate body; pregeniculatum) is an oval elevation on the lateral part of the posterior end of the thalamus, and is connected with the superior colliculus by the superior brachium. It is of a dark color, and presents a laminated arrangement consisting of alternate layers of gray and white substance. It receives numerous fibers from the optic tract, while other fibers of this tract pass over or through it into the pulvinar. Its cells are large and pigmented; their axons pass to the visual area in the occipital part of the cerebral cortex.
The superior colliculus, the pulvinar, and the lateral geniculate body receive many fibers from the optic tracts, and are therefore intimately connected with sight, constituting what are termed the lower visual centers. Extirpation of the eyes in a newly born animal entails an arrest of the development of these centers, but has no effect on the medial geniculate bodies or on the inferior colliculi. Moreover, the latter are well-developed in the mole, an animal in which the superior colliculi are rudimentary.
The trigonum habenulæ is a small depressed triangular area situated in front of the superior colliculus and on the lateral aspect of the posterior part of the tænia thalami. It contains a group of nerve cells termed the ganglion habenulæ. Fibers enter it from the stalk of the pineal body, and others, forming what is termed the habenular commissure, pass across the middle line to the corresponding ganglion of the opposite side. Most of its fibers are, however, directed downward and form a bundle, the fasciculus retroflexus of Meynert, which passes medial to the red nucleus, and, after decussating with the corresponding fasciculus of the opposite side, ends in the interpeduncular ganglion.
The pineal body (corpus pineale; epiphysis) is a small, conical, reddish-gray body which lies in the depression between the superior colliculi. It is placed beneath the splenium of the corpus callosum, but is separated from this by the tela chorioidea of the third ventricle, the lower layer of which envelops it. It measures about 8 mm. in length, and its base, directed forward, is attached by a stalk or peduncle of white substance. The stalk of the pineal body divides anteriorly into two laminæ, a dorsal and a ventral, separated from one another by the pineal recess of the third ventricle. The ventral lamina is continuous with the posterior commissure; the dorsal lamina is continuous with the habenular commissure and divides into two strands the medullary striæ, which run forward, one on either side, along the junction of the medial and upper surfaces of the thalamus to blend in front with the columns of the fornix.
The posterior commissure is a rounded band of white fibers crossing the middle line on the dorsal aspect of the upper end of the cerebral aqueduct. Its fibers acquire their medullary sheaths early, but their connections have not been definitely determined. Most of them have their origin in a nucleus, the nucleus of the posterior commissure (nucleus of Darkschewitsch), which lies in the central gray substance of the upper end of the cerebral aqueduct, in front of the nucleus of the oculomotor nerve. Some are probably derived from the posterior part of the thalamus and from the superior colliculus, while others are believed to be continued downward into the medial longitudinal fasciculus.
The Hypothalamus(Fig. 720) includes the subthalamic tegmental region and the structures forming the greater part of the floor of the third ventricle, viz., the corpora mammillaria, tuber cinereum, infundibulum, hypophysis, and optic chiasma.
The subthalamic tegmental region consists of the upward continuation of the tegmentum; it lies on the ventro-lateral aspect of the thalamus and separates it from the fibers of the internal capsule. The red nucleus and the substantia nigra are prolonged into its lower part; in front it is continuous with the substantia innominata of Meynert, medially with the gray substance of the floor of the third ventricle.
It consists from above downward of three strata: (1) stratum dorsale, directly applied to the under surface of the thalamus and consisting of fine longitudinal fibers; (2) zona incerta, a continuation forward of the formatio reticularis of the tegmentum; and (3) the corpus subthalamicum (nucleus of Luys), a brownish mass presenting a lenticular shape on transverse section, and situated on the dorsal aspect of the fibers of the base of the cerebral peduncle; it is encapsuled by a lamina of nerve fibers and contains numerous medium-sized nerve cells, the connections of which are as yet not fully determined.
The corpora mammillaria (corpus albicantia) are two round white masses, each about the size of a small pea, placed side by side below the gray substance of the floor of the third ventricle in front of the posterior perforated substance. They consist of white substance externally and of gray substance internally, the cells of the latter forming two nuclei, a medial of smaller and a lateral of larger cells. The white substance is mainly formed by the fibers of the columns of the fornix, which descend to the base of the brain and end partly in the corpora mammillaria. From the cells of the gray substance of each mammillary body two fasciculi arise: one, the thalamomammillary fasciculus (bundle of Vicq dAzyr), passes upward into the anterior nucleus of the thalamus; the other is directed downward into the tegmentum. Afferent fibers are believed to reach the corpus mammillare from the medial lemniscus and from the tegmentum.
FIG. 720 Median sagittal section of brain. The relations of the pia mater are indicated by the red color. (See enlarged image)
The tuber cinereum is a hollow eminence of gray substance situated between the corpora mammillaria behind, and the optic chiasma in front. Laterally it is continuous with the anterior perforated substances and anteriorly with a thin lamina, the lamina terminalis. From the under surface of the tuber cinereum a hollow conical process, the infundibulum, projects downward and forward and is attached to the posterior lobe of the hypophysis.
In the lateral part of the tuber cinereum is a nucleus of nerve cells, the basal optic nucleus of Meynert, while close to the cavity of the third ventricle are three additional nuclei. Between the tuber cinereum and the corpora mammillaria a small elevation, with a corresponding depression in the third ventricle, is sometimes seen. Retzius has named it the eminentia saccularis, and regards it as a representative of the saccus vasculosus found in this situation in some of the lower vertebrates.
The hypophysis (pituitary body) (Fig. 721) is a reddish-gray, somewhat oval mass, measuring about 12.5 mm. in its transverse, and about 8 mm. in its antero-posterior diameter. It is attached to the end of the infundibulum, and is situated in the fossa hypophyseos of the sphenoidal bone, where it is retained by a circular fold of dura mater, the diaphragma sella; this fold almost completely roofs in the fossa, leaving only a small central aperture through which the infundibulum passes.
FIG. 721 The hypophysis cerebri, in position. Shown in sagittal section. (See enlarged image)
Optic Chiasma (chiasma opticum; optic commissure).The optic chiasma is a flattened, somewhat quadrilateral band of fibers, situated at the junction of the floor and anterior wall of the third ventricle. Most of its fibers have their origins in the retina, and reach the chiasma through the optic nerves, which are continuous with its antero-lateral angles. In the chiasma, they undergo a partial decussation (Fig. 722); the fibers from the nasal half of the retina decussate and enter the optic tract of the opposite side, while the fibers from the temporal half of the retina do not undergo decussation, but pass back into the optic tract of the same side. Occupying the posterior part of the commissure, however, is a strand of fibers, the commissure of Gudden, which is not derived from the optic nerves; it forms a connecting link between the medial geniculate bodies.
Optic Tracts.The optic tracts are continued backward and lateralward from the postero-lateral angles of the optic chiasma. Each passes between the anterior perforated substance and the tuber cinereum, and, winding around the ventrolateral aspect of the cerebral peduncle, divides into a medial and a lateral root. The former comprises the fibers of Guddens commissure. The lateral root consists mainly of afferent fibers which arise in the retina and undergo partial decussation in the optic chiasma, as described; but it also contains a few fine efferent fibers which have their origins in the brain and their terminations in the retina. When traced backward, the afferent fibers of the lateral root are found to end in the lateral geniculate body and pulvinar of the thalamus, and in the superior colliculus; and these three structures constitute the lower visual centers. Fibers arise from the nerve cells in these centers and pass through the occipital part of the internal capsule, under the name of the optic radiations, to the cortex of the occipital lobe of the cerebrum, where the higher or cortical visual center is situated. Some of the fibers of the optic radiations take an opposite course, arising from the cells of the occipital cortex and passing to the lower visual centers. Some fibers are detached from the optic tract, and pass through the cerebral peduncle to the nucleus of the oculomotor nerve. These may be regarded as the afferent branches for the Sphincter pupillæ and Ciliaris muscles. Other fibers have been described as reaching the cerebellum through the superior peduncle; while others, again, are lost in the pons.
The Third Ventricle (ventriculus tertius) (Figs. 716,720).The third ventricle is a median cleft between the two thalami. Behind, it communicates with the fourth ventricle through the cerebral aqueduct, and in front with the lateral ventricles through the interventricular foramen. Somewhat triangular in shape, with the apex directed backward, it has a roof, a floor, an anterior and a posterior boundary and a pair of lateral walls.
FIG. 722 Scheme showing central connections of the optic nerves and optic tracts. (See enlarged image)
The roof(Fig. 723) is formed by a layer of epithelium, which stretches between the upper edges of the lateral walls of the cavity and is continuous with the epithelial lining of the ventricle. It is covered by and adherent to a fold of pia mater, named the tela chorioidea of the third ventricle, from the under surface of which a pair of vascular fringed processes, the choroid plexuses of the third ventricle, project downward, one on either side of the middle line, and invaginate the epithelial roof into the ventricular cavity.
The floor slopes downward and forward and is formed mainly by the structures which constitute the hypothalamus: from before backward these are: the optic chiasma, the tuber cinereum and infundibulum, and the corpora mammillaria. Behind the last, the floor is formed by the interpeduncular fossa and the tegmenta of the cerebral peduncles. The ventricle is prolonged downward as a funnel-shaped recess, the recessus infundibuli, into the infundibulum, and to the apex of the latter the hypophysis is attached.
The anterior boundary is constituted below by the lamina terminalis, a thin layer of gray substance stretching from the upper surface of the optic chiasma to the rostrum of the corpus callosum; above by the columns of the fornix and the anterior commissure. At the junction of the floor and anterior wall, immediately above the optic chiasma, the ventricle presents a small angular recess or diverticulum, the optic recess. Between the columns of the fornix, and above the anterior commissure, is a second recess termed the vulva. At the junction of the roof and anterior wall of the ventricle, and situated between the thalami behind and the columns of the fornix in front, is the interventricular foramen (foramen of Monro) through which the third communicates with the lateral ventricles.
FIG. 723 Coronal section of lateral and third ventricles. (Diagrammatic.) (See enlarged image)
The posterior boundary is constituted by the pineal body, the posterior commissure and the cerebral aqueduct. A small recess, the recessus pinealis, projects into the stalk of the pineal body, while in front of and above the pineal body is a second recess, the recessus suprapinealis, consisting of a diverticulum of the epithelium which forms the ventricular roof.
Each lateral wall consists of an upper portion formed by the medial surface of the anterior two-thirds of the thalamus, and a lower consisting of an upward continuation of the gray substance of the ventricular floor. These two parts correspond to the alar and basal laminæ respectively of the lateral wall of the fore-brain vesicle and are separated from each other by a furrow, the sulcus of Monro, which extends from the interventricular foramen to the cerebral aqueduct (pages 741 and 742). The lateral wall is limited above by the tænia thalami. The columns of the fornix curve downward in front of the interventricular foramen, and then run in the lateral walls of the ventricle, where, at first, they form distinct prominences, but subsequently are lost to sight. The lateral walls are joined to each other across the cavity of the ventricle by a band of gray matter, the massa intermedia (page 809).
Interpeduncular Fossa (Fig. 724).This is a somewhat lozenge-shaped area of the base of the brain, limited in front by the optic chiasma, behind by the antero-superior surface of the pons, antero-laterally by the converging optic tracts, and postero-laterally by the diverging cerebral peduncles. The structures contained in it have already been described; from behind forward, they are the posterior perforated substance, corpora mamillaria, tuber cinereum, infundibulum, and hypophysis.
The Telencephalon.The telencephalon includes: (1) the cerebral hemispheres with their cavities, the lateral ventricles; and (2) the pars optica hypothalami and the anterior portion of the third ventricle (already described under the diencephalon). As previously stated (see page 744), each cerebral hemisphere may be divided into three fundamental parts, viz., the rhinencephalon, the corpus striatum, and the neopallium. The rhinencephalon, associated with the sense of smell, is the oldest part of the telencephalon, and forms almost the whole of the hemisphere in some of the lower animals, e. g., fishes, amphibians, and reptiles. In man it is rudimentary, whereas the neopallium undergoes great development and forms the chief part of the hemisphere.
The Cerebral Hemispheres.The cerebral hemispheres constitute the largest part of the brain, and, when viewed together from above, assume the form of an ovoid mass broader behind than in front, the greatest transverse diameter corresponding with a line connecting the two parietal eminences. The hemispheres are separated medially by a deep cleft, named the longitudinal cerebral fissure, and each possesses a central cavity, the lateral ventricle.
The Longitudinal Cerebral Fissure (fissura cerebri longitudinalis; great longitudinal fissure) contains a sickle-shaped process of dura mater, the falx cerebri. It front and behind, the fissure extends from the upper to the under surfaces of the hemispheres and completely separates them, but its middle portion separates them for only about one-half of their vertical extent; for at this part they are connected across the middle line by a great central white commissure, the corpus callosum.
In a median sagittal section (Fig. 720) the cut corpus callosum presents the appearance of a broad, arched band. Its thick posterior end, termed the splenium, overlaps the mid-brain, but is separated from it by the tela chorioidea of the third ventricle and the pineal body. Its anterior curved end, termed the genu, gradually tapers into a thinner portion, the rostrum, which is continued downward and backward in front of the anterior commissure to join the lamina terminalis. Arching backward from immediately behind the anterior commissure to the under surface of the splenium is a second white band named the fornix: between this and the corpus callosum are the laminæ and cavity of the septum pellucidum.
The lateral surface is convex in adaptation to the concavity of the corresponding half of the vault of the cranium. The medial surface is flat and vertical, and is separated from that of the opposite hemisphere by the great longitudinal fissure and the falx cerebri. The inferior surface is of an irregular form, and may be divided into three areas: anterior, middle, and posterior. The anterior area, formed by the orbital surface of the frontal lobe, is concave, and rests on the roof of the orbit and nose; the middle area is convex, and consists of the under surface of the temporal lobe: it is adapted to the corresponding half of the middle cranial fossa. The posterior area is concave, directed medialward as well as downward, and is named the tentorial surface, since it rests upon the tentorium cerebelli, which intervenes between it and the upper surface of the cerebellum.
FIG. 725 Lateral surface of left cerebral hemisphere, viewed from above. (See enlarged image)
These three surfaces are separated from each other by the following borders: (a) supero-medial, between the lateral and medial surfaces; (b) infero-lateral, between the lateral and inferior surfaces; the anterior part of this border separating the lateral from the orbital surface, is known as the superciliary border; (c) medial occipital, separating the medial and tentorial surfaces; and (d) medial orbital, separating the orbital from the medial surface. The anterior end of the hemisphere is named the frontal pole; the posterior, the occipital pole; and the anterior end of the temporal lobe, the temporal pole. About 5 cm. in front of the occipital pole on the infero-lateral border is an indentation or notch, named the preoccipital notch.
The surfaces of the hemispheres are moulded into a number of irregular eminences, named gyri or convolutions, and separated by furrows termed fissures and sulci. The furrows are of two kinds, complete and incomplete. The former appear early in fetal life, are few in number, and are produced by infoldings of the entire thickness of the brain wall, and give rise to corresponding elevations in the interior of the ventricle. They comprise the hippocampal fissure, and parts of the calcarine and collateral fissures. The incomplete furrows are very numerous, and only indent the subjacent white substance, without producing any corresponding elevations in the ventricular cavity.
The gyri and their intervening fissures and the sulci are fairly constant in their arrangement; at the same time they vary within certain limits, not only in different individuals, but on the two hemispheres of the same brain. The convoluted condition of the surface permits of a great increase of the gray matter without the sacrifice of much additional space. The number and extent of the gyri, as well as the depth of the intervening furrows, appear to bear a direct relation to the intellectual powers of the individual.
Certain of the fissures and sulci are utilized for the purpose of dividing the hemisphere into lobes, and are therefore termed interlobular; included under this category are the lateral cerebral, parietoöccipital, calcarine, and collateral fissures, the central and cingulate sulci, and the sulcus circularis.
FIG. 726 Lateral surface of left cerebral hemisphere, viewed from the side. (See enlarged image)
The Lateral Cerebral Fissure (fissura cerebri lateralis [Sylvii]; fissure of Sylvius) (Fig. 726) is a well-marked cleft on the inferior and lateral surfaces of the hemisphere, and consists of a short stem which divides into three rami. The stem is situated on the base of the brain, and commences in a depression at the lateral angle of the anterior perforated substance. From this point it extends between the anterior part of the temporal lobe and the orbital surface of the frontal lobe, and reaches the lateral surface of the hemisphere. Here it divides into three rami: an anterior horizontal, an anterior ascending, and a posterior. The anterior horizontal ramus passes foward for about 2.5 cm. into the inferior frontal gyrus, while the anterior ascending ramus extends upward into the same convolution for about an equal distance. The posterior ramus is the longest; it runs backward and slightly upward for about 7 cm., and ends by an upward inflexion in the parietal lobe.
The Central Sulcus (sulcus centralis [Rolandi]; fissure of Rolando; central fissure) (Figs. 725,726) is situated about the middle of the lateral surface of the hemisphere, and begins in or near the longitudinal cerebral fissure, a little behind its mid-point. It runs sinuously downward and forward, and ends a little above the posterior ramus of the lateral fissure, and about 2.5 cm. behind the anterior ascending ramus of the same fissure. It described two chief curves: a superior genu with its concavity directed forward, and an inferior genu with its concavity directed backward. The central sulcus forms an angle opening forward of about 70° with the median plane.
The medial part of the parietoöccipital fissure (Fig. 727) runs downward and forward as a deep cleft on the medial surface of the hemisphere, and joins the calcarine fissure below and behind the posterior end of the corpus callosum. In most cases it contains a submerged gyrus.
The Calcarine Fissure (fissura calcarina) (Fig. 727) is on the medial surface of the hemisphere. It begins near the occipital pole in two converging rami, and runs forward to a point a little below the splenium of the corpus callosum, where it is joined at an acute angle by the medial part of the parietoöccipital fissure. The anterior part of this fissure gives rise to the prominence of the calcar avis in the posterior cornu of the lateral ventricle.
The Cingulate Sulcus (sulcus cinguli; callosomarginal fissure) (Fig. 727) is on the medial surface of the hemisphere; it begins below the anterior end of the corpus callosum and runs upward and forward nearly parallel to the rostrum of this body and, curving in front of the genu, is continued backward above the corpus callosum, and finally ascends to the supero-medial border of the hemisphere a short distance behind the upper end of the central sulcus. It separates the superior frontal from the cingulate gyrus.
The Collateral Fissure (fissura collateralis) (Fig. 727) is on the tentorial surface of the hemisphere and extends from near the occipital pole to within a short distance of the temporal pole. Behind, it lies below and lateral to the calcarine fissure, from which it is separated by the lingual gyrus; in front, it is situated between the hippocampal gyrus and the anterior part of the fusiform gyrus.
The Sulcus Circularis (circuminsular fissure) (Fig. 731) is on the lower and lateral surfaces of the hemisphere: it surrounds the insula and separates it from the frontal, parietal, and temporal lobes.
Lobes of the Hemispheres.By means of these fissures and sulci, assisted by certain arbitrary lines, each hemisphere is divided into the following lobes: the frontal, the parietal, the temporal, the occipital, the limbic, and the insula.
Frontal Lobe (lobus frontalis).On the lateral surface of the hemisphere this lobe extends from the frontal pole to the central sulcus, the latter separating it from the parietal lobe. Below, it is limited by the posterior ramus of the lateral fissure, which intervenes between it and the central lobe. On the medial surface, it is separated from the cingulate gyrus by the cingulate sulcus; and on the inferior surface, it is bounded behind by the stem of the lateral fissure.
FIG. 728 Principal fissures and lobes of the cerebrum viewed laterally. (See enlarged image)
The lateral surface of the frontal lobe (Fig. 726) is tranversed by three sulci which divide it into four gyri: the sulci are named the precentral, and the superior and inferior frontal; the gyri are the anterior central, and the superior, middle, and inferior frontal. The precentral sulcus runs parallel to the central sulcus, and is usually divided into an upper and a lower part; between it and the central sulcus is the anterior central gyrus. From the precentral sulcus, the superior and inferior frontal sulci run forward and downward, and divide the remainder of the lateral surface of the lobe into three parallel gyri, named, respectively the superior, middle, and inferior frontal gyri.
The anterior central gyrus (gyrus centralis anterior; ascending frontal convolution; precentral gyre) is bounded in front by the precentral sulcus, behind by the central sulcus; it extends from the supero-medial border of the hemisphere to the posterior ramus of the lateral fissure.
The superior frontal gyrus (gyrus frontalis superior; superfrontal gyre) is situated above the superior frontal sulcus and is continued on to the medial surface of the hemisphere. The portion on the lateral surface of the hemisphere is usually more or less completely subdivided into an upper and a lower part by an antero-posterior sulcus, the paramedial sulcus, which, however, is frequently interrupted by bridging gyri.
The middle frontal gyrus (gyrus frontalis medius; medifrontal gyre), between the superior and inferior frontal sulci, is continuous with the anterior orbital gyrus on the inferior surface of the hemisphere; it is frequently subdivided into two by a horizontal sulcus, the medial frontal sulcus of Eberstaller, which ends anteriorly in a wide bifurcation.
The inferior frontal gyrus (gyrus frontalis inferior; subfrontal gyre) lies below the inferior frontal sulcus, and extends forward from the lower part of the precentral sulcus; it is continuous with the lateral and posterior orbital gyri on the under surface of the lobe. It is subdivided by the anterior horizontal and ascending rami of the lateral fissure into three parts, viz., (1) the orbital part, below the anterior horizontal ramus of the fissure; (2) the triangular part (cap of Broca), between the ascending and horizontal rami; and (3) the basilar part, behind the anterior ascending ramus. The left inferior frontal gyrus is, as a rule, more highly developed than the right, and is named the gyrus of Broca, from the fact that Broca described it as the center for articulate speech.
The inferior or orbital surface of the frontal lobe is concave, and rests on the orbital plate of the frontal bone (Fig. 729). It is divided into four orbital gyri by a well-marked H-shaped orbital sulcus. These are named, from their position, the medial, anterior, lateral, and posterior orbital gyri. The medial orbital gyrus presents a well-marked antero-posterior sulcus, the olfactory sulcus, for the olfactory tract; the portion medial to this is named the straight gyrus, and is continuous with the superior frontal gyrus on the medial surface.
The medial surface of the frontal lobe is occupied by the medial part of the superior frontal gyrus (marginal gyrus) (Fig. 727). It lies between the cingulate sulcus and the supero-medial margin of the hemisphere. The posterior part of this gyrus is sometimes marked off by a vertical sulcus, and is distinguished as the paracentral lobule, because it is continuous with the anterior and posterior central gyri.
Parietal Lobe (lobus parietalis).The parietal lobe is separated from the frontal lobe by the central sulcus, but its boundaries below and behind are not so definite. Posteriorly, it is limited by the parietoöccipital fissure, and by a line carried across the hemisphere from the end of this fissure toward the preoccipital notch. Below, it is separated from the temporal lobe by the posterior ramus of the lateral fissure, and by a line carried backward from it to meet the line passing downward to the preoccipital notch.
The lateral surface of the parietal lobe (Fig. 726) is cleft by a well-marked furrow, the intraparietal sulcus of Turner, which consists of an oblique and a horizontal portion. The oblique part is named the postcentral sulcus, and commences below, about midway between the lower end of the central sulcus and the upturned end of the lateral fissure. It runs upward and backward, parallel to the central sulcus, and is sometimes divided into an upper and a lower ramus. It forms the hinder limit of the posterior central gyrus.
From about the middle of the postcentral sulcus, or from the upper end of its inferior ramus, the horizontal portion of the intraparietal sulcus is carried backward and slightly upward on the parietal lobe, and is prolonged, under the name of the occipital ramus, on to the occipital lobe, where it divides into two parts, which form nearly a right angle with the main stem and constitute the transverse occipital sulcus. The part of the parietal lobe above the horizontal portion of the intraparietal sulcus is named the superior parietal lobule; the part below, the inferior parietal lobule.
The posterior central gyrus (gyrus centralis posterior; ascending parietal convolution; postcentral gyre) extends from the longitudinal fissure above to the posterior ramus of the lateral fissure below. It lies parallel with the anterior central gyrus, with which it is connected below, and also, sometimes, above, the central sulcus.
The superior parietal lobule (lobulus parietalis superior) is bounded in front by the upper part of the postcentral sulcus, but is usually connected with the posterior central gyrus above the end of the sulcus; behind it is the lateral part of the parietoöccipital fissure, around the end of which it is joined to the occipital lobe by a curved gyrus, the arcus parietoöccipitalis; below, it is separated from the inferior parietal lobule by the horizontal portion of the intraparietal sulcus.
The inferior parietal lobule (lobulus parietalis inferior; subparietal district or lobule) lies below the horizontal portion of the intraparietal sulcus, and behind the lower part of the postcentral sulcus. It is divided from before backward into two gyri. One, the supramarginal, arches over the upturned end of the lateral fissure; it is continuous in front with the postcentral gyrus, and behind with the superior temporal gyrus. The second, the angular, arches over the posterior end of the superior temporal sulcus, behind which it is continuous with the middle temporal gyrus.
The medial surface of the parietal lobe (Fig. 727) is bounded behind by the medial part of the parietoöccipital fissure; in front, by the posterior end of the cingulate sulcus; and below, it is separated from the cingulate gyrus by the subparietal sulcus. It is of small size, and consists of a square-shaped convolution, which is termed the precuneus or quadrate lobe.
The lateral surface is limited in front by the lateral part of the parietoöccipital fissure, and by a line carried from the end of this fissure to the preoccipital notch; it is traversed by the transverse occipital and the lateral occipital sulci. The transverse occipital sulcus is continuous with the posterior end of the occipital ramus of the intraparietal sulcus, and runs across the upper part of the lobe, a short distance behind the parietoöccipital fissure. The lateral occipital sulcus extends from behind forward, and divides the lateral surface of the occipital lobe into a superior and an inferior gyrus, which are continuous in front with the parietal and temporal lobes.125
The medial surface of the occipital lobe is bounded in front by the medial part of the parietoöccipital fissure, and is traversed by the calcarine fissure, which subdivides it into the cuneus and the lingual gyrus. The cuneus is a wedge-shaped area between the calcarine fissure and the medial part of the parietoöccipital fissure. The lingual gyrus lies between the calcarine fissure and the posterior part of the collateral fissure; behind, it reaches the occipital pole; in front, it is continued on to the tentorial surface of the temporal lobe, and joins the hippocampal gyrus.
The tentorial surface of the occipital lobe is limited in front by an imaginary transverse line through the preoccipital notch, and consists of the posterior part of the fusiform gyrus (occipitotemporal convolution) and the lower part of the lingual gyrus, which are separated from each other by the posterior segment of the collateral fissure.
The superior surface forms the lower limit of the lateral fissure and overlaps the insula. On opening out the lateral fissure, three or four gyri will be seen springing from the depth of the hinder end of the fissure, and running obliquely forward and outward on the posterior part of the upper surface of the superior temporal gyrus; these are named the transverse temporal gyri (Heschl) (Fig. 730).
The lateral surface(Fig. 726) is bounded above by the posterior ramus of the lateral fissure, and by the imaginary line continued backward from it; below, it is limited by the infero-lateral border of the hemisphere. It is divided into superior, middle, and inferior gyri by the superior and middle temporal sulci. The superior temporal sulcus runs from before backward across the temporal lobe, some little distance below, but parallel with, the posterior ramus of the lateral fissure; and hence it is often termed the parallel sulcus. The middle temporal sulcus takes the same direction as the superior, but is situated at a lower level, and is usually subdivided into two or more parts. The superior temporal gyrus lies between the posterior ramus of the lateral fissure and the superior temporal sulcus, and is continuous behind with the supramarginal and angular gyri. The middle temporal gyrus is placed between the superior and middle temporal sulci, and is joined posteriorly with the angular gyrus. The inferior temporal gyrus is placed below the middle temporal sulcus, and is connected behind with the inferior occipital gyrus; it also extends around the infero-lateral border on to the inferior surface of the temporal lobe, where it is limited by the inferior sulcus.
The inferior surface is concave, and is continuous posteriorly with the tentorial surface of the occipital lobe. It is traversed by the inferior temporal sulcus, which extends from near the occipital pole behind, to within a short distance of the temporal pole in front, but is frequently subdivided by bridging gyri. Lateral to this fissure is the narrow tentorial part of the inferior temporal gyrus, and medial to it the fusiform gyrus, which extends from the occipital to the temporal pole; this gyrus is limited medially by the collateral fissure, which separates it from the lingual gyrus behind and from the hippocampal gyrus in front.
The Insula (island of Reil; central lobe) (Fig. 731) lies deeply in the lateral or Sylvian fissure, and can only be seen when the lips of that fissure are widely separated, since it is overlapped and hidden by the gyri which bound the fissure. These gyri are termed the opercula of the insula; they are separated from each other by the three rami of the lateral fissure, and are named the orbital, frontal, frontoparietal, and temporal opercula. The orbital operculum lies below the anterior horizontal ramus of the fissure, the frontal between this and the anterior ascending ramus, the parietal between the anterior ascending ramus and the upturned end of the posterior ramus, and the temporal below the posterior ramus. The frontal operculum is of small size in those cases where the anterior horizontal and ascending rami of the lateral fissure arise from a common stem. The insula is surrounded by a deep circular sulcus which separates it from the frontal, parietal, and temporal lobes. When the opercula have been removed, the insula is seen as a triangular eminence, the apex of which is directed toward the anterior perforated substance. It is divided into a larger anterior and a smaller posterior part by a deep sulcus, which runs backward and upward from the apex of the insula. The anterior part is subdivided by shallow sulci into three or four short gyri, while the posterior part is formed by one long gyrus, which is often bifurcated at its upper end. The cortical gray substance of the insula is continuous with that of the different opercula, while its deep surface corresponds with the lentiform nucleus of the corpus striatum.
FIG. 731 The insula of the left side, exposed by removing the opercula. (See enlarged image)
Limbic Lobe (Fig. 727).The term limbic lobe was introduced by Broca, and under it he included the cingulate and hippocampal gyri, which together arch around the corpus callosum and the hippocampal fissure. These he separated on the morphological ground that they are well-developed in animals possessing a keen sense of smell (osmatic animals), such as the dog and fox. They were thus regarded as a part of the rhinencephalon, but it is now recognized that they belong to the neopallium; the cingulate gyrus is therefore sometimes described as a part of the frontal lobe, and the hippocampal as a part of the temporal lobe.
The cingulate gyrus (gyrus cinguli; callosal convolution) is an arch-shaped convolution, lying in close relation to the superficial surface of the corpus callosum, from which it is separated by a slit-like fissure, the callosal fissure. It commences below the rostrum of the corpus callosum, curves around in front of the genu, extends along the upper surface of the body, and finally turns downward behind the splenium, where it is connected by a narrow isthmus with the hippocampal gyrus. It is separated from the medial part of the superior frontal gyrus by the cingulate sulcus, and from the precuneus by the subparietal sulcus.
The hippocampal gyrus (gyrus hippocampi) is bounded above by the hippocampal fissure, and below by the anterior part of the collateral fissure. Behind, it is continuous superiorly, through the isthmus, with the cingulate gyrus and inferiorly with the lingual gyrus. Running in the substance of the cingulate and hippocampal gyri, and connecting them together, is a tract of arched fibers, named the cingulum (page 843). The anterior extremity of the hippocampal gyrus is recurved in the form of a hook (uncus), which is separated from the apex of the temporal lobe by a slight fissure, the incisura temporalis. Although superficially continuous with the hippocampal gyrus, the uncus forms morphologically a part of the rhinencephalon.
The Hippocampal Fissure (fissura hippocampi; dentate fissure) begins immediately behind the splenium of the corpus callosum, and runs forward between the hippocampal and dentate gyri to end in the uncus. It is a complete fissure (page 819), and gives rise to the prominence of the hippocampus in the inferior cornu of the lateral ventricle.
Rhinencephalon (Fig. 732).The rhinencephalon comprises the olfactory lobe, the uncus, the subcallosal and supracallosal gyri, the fascia dentata hippocampi, the septum pellucidum, the fornix, and the hippocampus.
1. The Olfactory Lobe (lobus olfactorius) is situated under the inferior or orbital surface of the frontal lobe. In many vertebrates it constitutes a well-marked portion of the hemisphere and contains an extension of the lateral ventricle; but in man and some other mammals it is rudimentary. It consists of the olfactory bulb and tract, the olfactory trigone, the parolfactory area of Broca, and the anterior perforated substance.
(a) The olfactory bulb (bulbus olfactorius) is an oval, reddish-gray mass which rests on the cribriform plate of the ethmoid and forms the anterior expanded extremity of the olfactory tract. Its under surface receives the olfactory nerves, which pass upward through the cribriform plate from the olfactory region of the nasal cavity. Its minute structure is described on page 848.
(b) The olfactory tract (tractus olfactorius) is a narrow white band, triangular on coronal section, the apex being directed upward. It lies in the olfactory sulcus on the inferior surface of the frontal lobe, and divides posteriorly into two striæ, a medial and a lateral. The lateral stria is directed across the lateral part of the anterior perforated substance and then bends abruptly medialward toward the uncus of the hippocampal gyrus. The medial stria turns medialward behind the parolfactory area and ends in the subcallosal gyrus; in some cases a small intermediate stria is seen running backward to the anterior perforated substance.
(c) The olfactory trigone (trigonum olfactorium) is a small triangular area in front of the anterior perforated substance. Its apex, directed forward, occupies the posterior part of the olfactory sulcus, and is brought into view by throwing back the olfactory tract.
(d) The parolfactory area of Broca (area parolfactoria) is a small triangular field on the medial surface of the hemisphere in front of the subcallosal gyrus, from which it is separated by the posterior parolfactory sulcus; it is continuous below with the olfactory trigone, and above and in front with the cingulate gyrus; it is limited anteriorly by the anterior parolfactory sulcus.
(e) The anterior perforated substance (substantia perforata anterior) is an irregularly quadrilateral area in front of the optic tract and behind the olfactory trigone, from which it is separated by the fissure prima; medially and in front it is continuous with the subcallosal gyrus; laterally it is bounded by the lateral stria of the olfactory tract and is continued into the uncus. Its gray substance is confluent above with that of the corpus striatum, and is perforated anteriorly by numerous small bloodvessels.
3. The Subcallosal, Supracallosal, and Dentate Gyri form a rudimentary arch-shaped lamina of gray substance extending over the corpus callosum and above the hippocampal gyrus from the anterior perforated substance to the uncus.
(a) The subcallosal gyrus (gyrus subcallosus; peduncle of the corpus callosum) is a narrow lamina on the medial surface of the hemisphere in front of the lamina terminalis, behind the parolfactory area, and below the rostrum of the corpus callosum. It is continuous around the genu of the corpus callosum with the supracallosal gyrus.
(b) The supracallosal gyrus (indusium griseum; gyrus epicallosus) consists of a thin layer of gray substance in contact with the upper surface of the corpus callosum and continuous laterally with the gray substance of the cingulate gyrus. It contains two longitudinally directed strands of fibers termed respectively the medial and lateral longitudinal striæ. The supracallosal gyrus is prolonged around the splenium of the corpus callosum as a delicate lamina, the fasciola cinerea, which is continuous below with the fascia dentata hippocampi.
(c) The fascia dentata hippocampi (gyrus dentatus) is a narrow band extending downward and forward above the hippocampal gyrus but separated from it by the hippocampal fissure; its free margin is notched and overlapped by the fimbriathe fimbriodentate fissure intervening. Anteriorly it is continued into the notch of the uncus, where it forms a sharp bend and is then prolonged as a delicate band, the band of Giacomini, over the uncus, on the lateral surface of which it is lost.
Interior of the Cerebral Hemispheres.If the upper part of either hemisphere be removed, at a level about 1.25 cm. above the corpus callosum, the central white substance will be exposed as an oval-shaped area, the centrum ovale minus, surrounded by a narrow convoluted margin of gray substance, and studded with numerous minute red dots (puncta vasculosa), produced by the escape of blood from divided bloodvessels. If the remaining portions of the hemispheres be slightly drawn apart a broad band of white substance, the corpus callosum, will be observed, connecting them at the bottom of the longitudinal fissure; the margins of the hemispheres which overlap the corpus callosum are called the labia cerebri. Each labrium is part of the cingulate gyrus already described; and the slit-like interval between it and the upper surface of the corpus callosum is termed the callosal fissure(Fig. 727). If the hemispheres be sliced off to a level with the upper surface of the corpus callosum, the white substance of that structure will be seen connecting the two hemispheres. The large expanse of medullary matter now exposed, surrounded by the convoluted margin of gray substance, is called the centrum ovale majus.
The Corpus Callosum(Fig. 733) is the great transverse commissure which unites the cerebral hemispheres and roofs in the lateral ventricles. A good conception of its position and size is obtained by examining a median sagittal section of the brain (Fig. 720), when it is seen to form an arched structure about 10 cm. long. Its anterior end is about 4 cm. from the frontal pole, and its posterior end about 6 cm. from the occipital pole of the hemisphere.
The anterior end is named the genu, and is bent downward and backward in front of the septum pellucidum; diminishing rapidly in thickness, it is prolonged backward under the name of the rostrum, which is connected below with the lamina terminalis. The anterior cerebral arteries are in contact with the under surface of the rostrum; they then arch over the front of the genu, and are carried backward above the body of the corpus callosum.
The posterior end is termed the splenium and constitutes the thickest part of the corpus callosum. It overlaps the tela chorioidea of the third ventricle and the mid-brain, and ends in a thick, convex, free border. A sagittal section of the splenium shows that the posterior end of the corpus callosum is acutely bent forward, the upper and lower parts being applied to each other.
The superior surface is convex from before backward, and is about 2.5 cm. wide. Its medial part forms the bottom of the longitudinal fissure, and is in contact posteriorly with the lower border of the falx cerebri. Laterally it is overlapped by the cingulate gyrus, but is separated from it by the slit-like callosal fissure. It is traversed by numerous transverse ridges and furrows, and is covered by a thin layer of gray matter, the supracallosal gyrus, which exhibits on either side of the middle line the medial and lateral longitudinal striæ, already described (page 827).
The inferior surface is concave, and forms on either side of the middle line the roof of the lateral ventricle. Medially, this surface is attached in front to the septum pellucidum; behind this it is fused with the upper surface of the body of the fornix, while the splenium is in contact with the tela chorioidea.
On either side, the fibers of the corpus callosum radiate in the white substance and pass to the various parts of the cerebral cortex; those curving forward from the genu into the frontal lobe constitute the forceps anterior, and those curving backward into the occipital lobe, the forceps posterior. Between these two parts is the main body of the fibers which constitute the tapetum and extend laterally on either side into the temporal lobe, and cover in the central part of the lateral ventricle.
FIG. 734 Scheme showing relations of the ventricles to the surface of the brain. (See enlarged image)
The Lateral Ventricles (ventriculus lateralis) (Fig. 734).The two lateral ventricles are irregular cavities situated in the lower and medial parts of the cerebral hemispheres, one on either side of the middle line. They are separated from each other by a median vertical partition, the septum pellucidum, but communicate with the third ventricle and indirectly with each other through the interventricular foramen. They are lined by a thin, diaphanous membrane, the ependyma, covered by ciliated epithelium, and contain cerebrospinal fluid, which, even in health, may be secreted in considerable amount. Each lateral ventricle consists of a central part or body, and three prolongations from it, termed cornua (Figs. 735,736).
The central part (pars centralis ventriculi lateralis; cella) (Fig. 737) of the lateral ventricle extends from the interventricular foramen to the splenium of the corpus callosum. It is an irregularly curved cavity, triangular on transverse section, with a roof, a floor, and a medial wall. The roof is formed by the under surface of the corpus callosum; the floor by the following parts, enumerated in their order of position, from before backward: the caudate nucleus of the corpus striatum, the stria terminalis and the terminal vein, the lateral portion of the upper surface of the thalamus, the choroid plexus, and the lateral part of the fornix; the medial wall is the posterior part of the septum pellucidum, which separates it from the opposite ventricle.
FIG. 735 Drawing of a cast of the ventricular cavities, viewed from above. (Retzius.) (See enlarged image)
FIG. 736 Drawing of a cast of the ventricular cavities, viewed from the side. (Retzius.) (See enlarged image)
The anterior cornu (cornu anterius; anterior horn; precornu) (Fig. 736) passes forward and lateralward, with a slight inclination downward, from the interventricular foramen into the frontal lobe, curving around the anterior end of the caudate nucleus. Its floor is formed by the upper surface of the reflected portion of the corpus callosum, the rostrum. It is bounded medially by the anterior portion of the septum pellucidum, and laterally by the head of the caudate nucleus. Its apex reaches the posterior surface of the genu of the corpus callosum.
The posterior cornu (cornu posterius; postcornu) (Figs. 737,788) passes into the occipital lobe, its direction being backward and lateralward, and then medialward. Its roof is formed by the fibers of the corpus callosum passing to the temporal and occipital lobes. On its medial wall is a longitudinal eminence, the calcar avis (hippocampus minor), which is an involution of the ventricular wall produced by the calcarine fissure. Above this the forceps posterior of the corpus callosum, sweeping around to enter the occipital lobe, causes another projection, termed the bulb of the posterior cornu. The calcar avis and bulb of the posterior cornu are extremely variable in their degree of development; in some cases they are ill-defined, in others prominent.
FIG. 737 Central part and anterior and posterior cornua of lateral ventricles exposed from above. (See enlarged image)
The inferior cornu (cornu inferior; descending horn; middle horn; medicornu) (Fig. 739), the largest of the three, traverses the temporal lobe of the brain, forming in its course a curve around the posterior end of the thalamus. It passes at first backward, lateralward, and downward, and then curves forward to within 2.5 cm. of the apex of the temporal lobe, its direction being fairly well indicated on the surface of the brain by that of the superior temporal sulcus. Its roof is formed chiefly by the inferior surface of the tapetum of the corpus callosum, but the tail of the caudate nucleus and the stria terminalis also extend forward in the roof of the inferior cornu to its extremity; the tail of the caudate nucleus joins the putamen. Its floor presents the following parts: the hippocampus, the fimbria hippocampi, the collateral eminence, and the choroid plexus. When the choroid plexus is removed, a cleft-like opening is left along the medial wall of the inferior cornu; this cleft constitutes the lower part of the choroidal fissure.
FIG. 738 Coronal section through posterior cornua of lateral ventricle. (See enlarged image)
FIG. 739 Posterior and inferior cornua of left lateral ventricle exposed from the side. (See enlarged image)
The hippocampus (hippocampus major) (Figs. 739,740) is a curved eminence, about 5 cm. long, which extends throughout the entire length of the floor of the inferior cornu. Its lower end is enlarged, and presents two or three rounded elevations or digitations which give it a paw-like appearance, and hence it is named the pes hippocampi. If a transverse section be made through the hippocampus, it will be seen that this eminence is produced by the folding of the wall of the hemisphere to form the hippocampal fissure. The main mass of the hippocampus consists of gray substance, but on its ventricular surface is a thin white layer, the alveus, which is continuous with the fimbria hippocampi.
The collateral eminence (eminentia collateralis) (Fig. 740) is an elongated swelling lying lateral to and parallel with the hippocampus. It corresponds with the middle part of the collateral fissure, and its size depends on the depth and direction of this fissure. It is continuous behind with a flattened triangular area, the trigonum collaterale, situated between the posterior and inferior cornua.
FIG. 741 Two views of a model of the striatum: A, lateral aspect; B, mesal aspect. (See enlarged image)
The corpus striatum has received its name from the striped appearance which a section of its anterior part presents, in consequence of diverging white fibers being mixed with the gray substance which forms its chief mass. A part of the corpus striatum is imbedded in the white substance of the hemisphere, and is therefore external to the ventricle; it is termed the extraventricular portion, or the lentiform nucleus; the remainder, however, projects into the ventricle, and is named the intraventricular portion, or the caudate nucleus(Fig. 737).
The caudate nucleus (nucleus caudatus; caudatum) (Figs. 741,742) is a pear-shaped, highly arched gray mass; its broad extremity, or head, is directed forward into the anterior cornu of the lateral ventricle, and is continuous with the anterior perforated substance and with the anterior end of the lentiform nucleus; its narrow end, or tail, is directed backward on the lateral side of the thalamus, from which it is separated by the stria terminalis and the terminal vein. It is then continued downward into the roof of the inferior cornu, and ends in the putamen near the apex of the temporal lobe. It is covered by the lining of the ventricle, and crossed by some veins of considerable size. It is separated from the lentiform nucleus, in the greater part of its extent, by a thick lamina of white substance, called the internal capsule, but the two portions of the corpus striatum are united in front (Figs. 743,744).
The lentiform nucleus (nucleus lentiformis; lenticular nucleus; lenticula) (Fig. 741) is lateral to the caudate nucleus and thalamus, and is seen only in sections of the hemisphere. When divided horizontally, it exhibits, to some extent, the appearance of a biconvex lens (Fig. 742), while a coronal section of its central part presents a somewhat triangular outline. It is shorter than the caudate nucleus and does not extend as far forward. It is bounded laterally by a lamina of white substance called the external capsule, and lateral to this is a thin layer of gray substance termed the claustrum. Its anterior end is continuous with the lower part of the head of the caudate nucleus and with the anterior perforated substance.
In a coronal section through the middle of the lentiform nucleus, two medullary laminæ are seen dividing it into three parts. The lateral and largest part is of a reddish color, and is known as the putamen, while the medial and intermediate are of a yellowish tint, and together constitute the globus pallidus; all three are marked by fine radiating white fibers, which are most distinct in the putamen (Fig. 744).
The gray substance of the corpus striatum is traversed by nerve fibers, some of which originate in it. The cells are multipolar, both large and small; those of the lentiform nucleus contain yellow pigment. The caudate and lentiform nuclei are not only directly continuous with each other anteriorly, but are connected to each other by numerous fibers. The corpus striatum is also connected: (1) to the cerebral cortex, by what are termed the corticostriate fibers; (2) to the thalamus, by fibers which pass through the internal capsule, and by a strand named the ansa lentiformis; (3) to the cerebral peduncle, by fibers which leave the lower aspect of the caudate and lentiform nuclei.
FIG. 743 Coronal section through anterior cornua of lateral ventricles. (See enlarged image)
The claustrum (Figs. 742,744) is a thin layer of gray substance, situated on the lateral surface of the external capsule. Its transverse section is triangular, with the apex directed upward. Its medial surface, contiguous to the external capsule, is smooth, but its lateral surface presents ridges and furrows corresponding with the gyri and sulci of the insula, with which it is in close relationship. The claustrum is regarded as a detached portion of the gray substance of the insula, from which it is separated by a layer of white fibers, the capsula extrema (band of Baillarger). Its cells are small and spindle-shaped, and contain yellow pigment; they are similar to those of the deepest layer of the cortex.
The nucleus amygdalæ (amygdala) (Fig. 741), is an ovoid gray mass, situated at the lower end of the roof of the inferior cornu. It is merely a localized thickening of the gray cortex, continuous with that of the uncus; in front it is continuous with the putamen, behind with the stria terminalis and the tail of the caudate nucleus.
The internal capsule (capsula interna) (Figs. 745,746) is a flattened band of white fibers, between the lentiform nucleus on the lateral side and the caudate nucleus and thalamus on the medial side. In horizontal section (Figs. 742) it is seen to be somewhat abruptly curved, with its convexity inward; the prominence of the curve is called the genu, and projects between the caudate nucleus and the thalamus. The portion in front of the genu is termed the frontal part, and separates the lentiform from the caudate nucleus; the portion behind the genu is the occipital part, and separates the lentiform nucleus from the thalamus.
The frontal part of the internal capsule contains: (1) fibers running from the thalamus to the frontal lobe; (2) fibers connecting the lentiform and caudate nuclei; (3) fibers connecting the cortex with the corpus striatum; and (4) fibers passing from the frontal lobe through the medial fifth of the base of the cerebral peduncle to the nuclei pontis. The fibers in the region of the genu are named the geniculate fibers; they originate in the motor part of the cerebral cortex, and, after passing downward through the base of the cerebral peduncle with the cerebrospinal fibers, undergo decussation and end in the motor nuclei of the cranial nerves of the opposite side. The anterior two-thirds of the occipital part of the internal capsule contains the cerebrospinal fibers, which arise in the motor area of the cerebral cortex and, passing downward through the middle three-fifths of the base of the cerebral peduncle, are continued into the pyramids of the medulla oblongata. The posterior third of the occipital part contains: (1) sensory fibers, largely derived from the thalamus, though some may be continued upward from the medial lemniscus; (2) the fibers of optic radiation, from the lower visual centers to the cortex of the occipital lobe; (3) acoustic fibers, from the lateral lemniscus to the temporal lobe; and (4) fibers which pass from the occipital and temporal lobes to the nuclei pontis.
The fibers of the internal capsule radiate widely as they pass to and from the various parts of the cerebral cortex, forming the corona radiata(Fig. 745) and intermingling with the fibers of the corpus callosum.
The external capsule (capsula externa) (Fig. 742) is a lamina of white substance, situated lateral to the lentiform nucleus, between it and the claustrum, and continuous with the internal capsule below and behind the lentiform nucleus. It probably contains fibers derived from the thalamus, the anterior commissure, and the subthalamic region.
FIG. 745 Dissection showing the course of the cerebrospinal fibers. (E. B. Jamieson.) (See enlarged image)
The substantia innominata of Meynert is a stratum consisting partly of gray and partly of white substance, which lies below the anterior part of the thalamus and lentiform nucleus. It consists of three layers, superior, middle, and inferior. The superior layer is named the ansa lentiformis, and its fibers, derived from the medullary lamina of the lentiform nucleus, pass medially to end in the thalamus and subthalamic region, while others are said to end in the tegmentum and red nucleus. The middle layer consists of nerve cells and nerve fibers; fibers enter it from the parietal lobe through the external capsule, while others are said to connect it with the medial longitudinal fasciculus. The inferior layer forms the main part of the inferior stalk of the thalamus, and connects this body with the temporal lobe and the insula.
The stria terminalis (tænia semicircularis) is a narrow band of white substance situated in the depression between the caudate nucleus and the thalamus. Anteriorly, its fibers are partly continued into the column of the fornix; some, however, pass over the anterior commissure to the gray substance between the caudate nucleus and septum pellucidum, while others are said to enter the caudate nucleus. Posteriorly, it is continued into the roof of the inferior cornu of the lateral ventricle, at the extremity of which it enters the nucleus amygdalæ. Superficial to it is a large vein, the terminal vein (vein of the corpus striatum), which receives numerous tributaries from the corpus striatum and thalamus; it runs forward to the interventricular foramen and there joins with the vein of the choroid plexus to form the corresponding internal cerebral vein. On the surface of the terminal vein is a narrow white band, named the lamina affixa.
The Fornix (Figs. 720,747,748) is a longitudinal, arch-shaped lamella of white substance, situated below the corpus callosum, and continuous with it behind, but separated from it in front by the septum pellucidum. It may be described as consisting of two symmetrical bands, one for either hemisphere. The two portions are not united to each other in front and behind, but their central parts are joined together in the middle line. The anterior parts are called the columns of the fornix; the intermediate united portions, the body; and the posterior parts, the crura.
FIG. 746 Diagram of the tracts in the internal capsule. Motor tract red. The sensory tract (blue) is not direct, but formed of neurons receiving impulses from below in the thalamus and transmitting them to the cortex. The optic radiation (occipitothalamic) is shown in violet. (See enlarged image)
The body (corpus fornicis) of the fornix is triangular, narrow in front, and broad behind. The medial part of its upper surface is connected to the septum pellucidum in front and to the corpus callosum behind. The lateral portion of this surface forms part of the floor of the lateral ventricle, and is covered by the ventricular epithelium. Its lateral edge overlaps the choroid plexus, and is continuous with the epithelial covering of this structure. The under surface rests upon the tela chorioidea of the third ventricle, which separates it from the epithelial roof of that cavity, and from the medial portions of the upper surfaces of the thalami. Below, the lateral portions of the body of the fornix are joined by a thin triangular lamina, named the psalterium (lyra). This lamina contains some transverse fibers which connect the two hippocampi across the middle line and constitute the hippocampal commissure. Between the psalterium and the corpus callosum a horizontal cleft, the so-called ventricle of the fornix (ventricle of Verga), is sometimes found.
The columns (columna fornicis; anterior pillars; fornicolumns) of the fornix arch downward in front of the interventricular foramen and behind the anterior commissure, and each descends through the gray substance in the lateral wall of the third ventricle to the base of the brain, where it ends in the corpus mammillare. From the cells of the corpus mammillare the thalamomammillary fasciculus (bundle of Vicq dAzyr) takes origin and is prolonged into the anterior nucleus of the thalamus. The column of the fornix and the thalamomammillary fasciculus together form a loop resembling the figure 8, but the continuity of the loop is broken in the corpus mammillare. The column of the fornix is joined by the stria medullaris of the pineal body and by the superficial fibers of the stria terminalis, and is said to receive also fibers from the septum pellucidum. Zuckerkandl describes an olfactory fasciculus which becomes detached from the main portion of the column of the fornix, and passes downward in front of the anterior commissure to the base of the brain, where it divides into two bundles, one joining the medial stria of the olfactory tract; the other joins the subcallosal gyrus, and through it reaches the hippocampal gyrus.
FIG. 748 The fornix and corpus callosum from below. (From a specimen in the Department of Human Anatomy of the University of Oxford.) (See enlarged image)
The crura (crus fornicis; posterior pillars) of the fornix are prolonged backward from the body. They are flattened bands, and at their commencement are intimately connected with the under surface of the corpus callosum. Diverging from one another, each curves around the posterior end of the thalamus, and passes downward and forward into the inferior cornu of the lateral ventricle (Fig. 750). Here it lies along the concavity of the hippocampus, on the surface of which some of its fibers are spread out to form the alveus, while the remainder are continued as a narrow white band, the fimbria hippocampi, which is prolonged into the uncus of the hippocampal gyrus. The inner edge of the fimbria overlaps the fascia dentata hippocampi (dentate gyrus) (page 827), from which it is separated by the fimbriodentate fissure; from its lateral margin, which is thin and ragged, the ventricular epithelium is reflected over the choroid plexus as the latter projects into the chorioidal fissure.
Interventricular Foramen (foramen of Monro).Between the columns of the fornix and the anterior ends of the thalami, an oval aperture is present on either side: this is the interventricular foramen, and through it the lateral ventricles communicate with the third ventricle. Behind the epithelial lining of the foramen the choroid plexuses of the lateral ventricles are joined across the middle line.
The Anterior Commissure (precommissure) is a bundle of white fibers, connecting the two cerebral hemispheres across the middle line, and placed in front of the columns of the fornix. On sagittal section it is oval in shape, its long diameter being vertical and measuring about 5 mm. Its fibers can be traced lateralward and backward on either side beneath the corpus striatum into the substance of the temporal lobe. It serves in this way to connect the two temporal lobes, but it also contains decussating fibers from the olfactory tracts.
The Septum Pellucidum (septum lucidum) (Fig. 720) is a thin, vertically placed partition consisting of two laminæ, separated in the greater part of their extent by a narrow chink or interval, the cavity of the septum pellucidum. It is attached, above, to the under surface of the corpus callosum; below, to the anterior part of the fornix behind, and the reflected portion of the corpus callosum in front. It is triangular in form, broad in front and narrow behind; its inferior angle corresponds with the upper part of the anterior commissure. The lateral surface of each lamina is directed toward the body and anterior cornu of the lateral ventricle, and is covered by the ependyma of that cavity.
The cavity of the septum pellucidum (cavum septi pellucidi; pseudocele; fifth ventricle) is generally regarded as part of the longitudinal cerebral fissure, which has become shut off by the union of the hemispheres in the formation of the corpus callosum above and the fornix below. Each half of the septum therefore forms part of the medial wall of the hemisphere, and consists of a medial layer of gray substance, derived from that of the cortex, and a lateral layer of white substance continuous with that of the cerebral hemispheres. This cavity is not developed from the cavity of the cerebral vesicles, and never communicates with the ventricles of the brain.
The Choroid Plexus of the Lateral Ventricle (plexus chorioideus ventriculus lateralis; paraplexus) (Fig. 750) is a highly vascular, fringe-like process of pia mater, which projects into the ventricular cavity. The plexus, however, is everywhere covered by a layer of epithelium continuous with the epithelial lining of the ventricle. It extends from the interventricular foramen, where it is joined with the plexus of the opposite ventricle, to the end of the inferior cornu. The part in relation to the body of the ventricle forms the vascular fringed margin of a triangular process of pia mater, named the tela chorioidea of the third ventricle, and projects from under cover of the lateral edge of the fornix. It lies upon the upper surface of the thalamus, from which the epithelium is reflected over the plexus on to the edge of the fornix (Fig. 723). The portion in relation to the inferior cornu lies in the concavity of the hippocampus and overlaps the fimbria hippocampi: from the lateral edge of the fimbria the epithelium is reflected over the plexus on to the roof of the cornu (Fig. 749). It consists of minute and highly vascular villous processes, each with an afferent and an efferent vessel. The arteries of the plexus are: (a) the anterior choroidal, a branch of the internal carotid, which enters the plexus at the end of the inferior cornu; and (b) the posterior choroidal, one or two small branches of the posterior cerebral, which pass forward under the splenium. The veins of the choroid plexus unite to form a tortuous vein, which courses from behind forward to the interventricular foramen and there joins with the terminal vein to form the corresponding internal cerebral vein.
FIG. 749 Coronal section of inferior horn of lateral ventricle. (Diagrammatic.) (See enlarged image)
When the choroid plexus is pulled away, the continuity between its epithelial covering and the epithelial lining of the ventricle is severed, and a cleft-like space is produced. This is named the choroidal fissure; like the plexus, it extends from the interventricular foramen to the end of the inferior cornu. The upper part of the fissure, i.e., the part nearest the interventricular foramen is situated between the lateral edge of the fornix and the upper surface of the thalamus; farther back at the beginning of the inferior cornu it is between the commencement of the fimbria hippocampi and the posterior end of the thalamus, while in the inferior cornu it lies between the fimbria in the floor and the stria terminalis in the roof of the cornu.
The tela chorioidea of the third ventricle (tela chorioidea ventriculi tertii; velum interpositum) (Fig. 750) is a double fold of pia mater, triangular in shape, which lies beneath the fornix. The lateral portions of its lower surface rest upon the thalami, while its medial portion is in contact with the epithelial roof of the third ventricle. Its apex is situated at the interventricular foramen; its base corresponds with the splenium of the corpus callosum, and occupies the interval between that structure above and the corpora quadrigemina and pineal body below. This interval, together with the lower portions of the choroidal fissures, is sometimes spoken of as the transverse fissure of the brain. At its base the two layers of the velum separate from each other, and are continuous with the pia mater investing the brain in this region. Its lateral margins are modified to form the highly vascular choroid plexuses of the lateral ventricles. It is supplied by the anterior and posterior choroidal arteries already described, The veins of the tela chorioidea are named the internal cerebral veins (venæ Galeni); they are two in number, and run backward between its layers, each being formed at the interventricular foramen by the union of the terminal vein with the choroidal vein. The internal cerebral veins unite posteriorly in a single trunk, the great cerebral vein (vena magna Galeni), which passes backward beneath the splenium and ends in the straight sinus.
FIG. 750 Tela chorioidea of the third ventricle, and the choroid plexus of the left lateral ventricle, exposed from above. (See enlarged image)
Structure of the Cerebral Hemispheres.The cerebral hemispheres are composed of gray and white substance: the former covers their surface, and is termed the cortex; the latter occupies the interior of the hemispheres.
The white substance consists of medullated fibers, varying in size, and arranged in bundles separated by neuroglia. They may be divided, according to their course and connections, into three distinct systems. (1) Projection fibers connect the hemisphere with the lower parts of the brain and with the medulla spinalis. (2) Transverse or commissural fibers unite the two hemispheres. (3) Association fibers connect different structures in the same hemisphere; these are, in many instances, collateral branches of the projection fibers, but others are the axons of independent cells.
1. The projection fibers consist of efferent and afferent fibers uniting the cortex with the lower parts of the brain and with the medulla spinalis. The principal efferent strands are: (1) the motor tract, occupying the genu and anterior two-thirds of the occipital part of the internal capsule, and consisting of (a) the geniculate fibers, which decussate and end in the motor nuclei of the cranial nerves of the opposite side; and (b) the cerebrospinal fibers, which are prolonged through the pyramid of the medulla oblongata into the medulla spinalis: (2) the corticopontine fibers, ending in the nuclei pontis. The chief afferent fibers are: (1) those of the lemniscus which are not interrupted in the thalamus; (2) those of the superior cerebellar peduncle which are not interrupted in the red nucleus and thalamus; (3) numerous fibers arising within the thalamus, and passing through its stalks to the different parts of the cortex (page 810); (4) optic and acoustic fibers, the former passing to the occipital, the latter to the temporal lobe.
2. The transverse or commissural fibers connect the two hemispheres. They include: (a) the transverse fibers of the corpus callosum, (b) the anterior commissure, (c) the posterior commissure, and (d) the lyra or hippocampal commissure; they have already been described.
FIG. 751 Diagram showing principal systems of association fibers in the cerebrum. (See enlarged image)
3. The association fibers(Fig. 751) unite different parts of the same hemisphere, and are of two kinds: (1) those connecting adjacent gyri, short association fibers; (2) those passing between more distant parts, long association fibers.
The long association fibers include the following: (a) the uncinate fasciculus; (b) the cingulum; (c) the superior longitudinal fasciculus; (d) the inferior longitudinal fasciculus; (e) the perpendicular fasciculus; (f) the occipitofrontal fasciculus; and (g) the fornix.
(b) The cingulum is a band of white matter contained within the cingulate gyrus. Beginning in front at the anterior perforated substance, it passes forward and upward parallel with the rostrum, winds around the genu, runs backward above the corpus callosum, turns around the splenium, and ends in the hippocampal gyrus.
(c) The superior longitudinal fasciculus passes backward from the frontal lobe above the lentiform nucleus and insula; some of its fibers end in the occipital lobe, and others curve downward and forward into the temporal lobe.
(f) The occipitofrontal fasciculus passes backward from the frontal lobe, along the lateral border of the caudate nucleus, and on the mesial aspect of the corona radiata; its fibers radiate in a fan-like manner and pass into the occipital and temporal lobes lateral to the posterior and inferior cornua. Déjerine regards the fibers of the tapetum as being derived from this fasciculus, and not from the corpus callosum.
(g) The fornix connects the hippocampal gyrus with the corpus mammillare and, by means of the thalamomammillary fasciculus, with the thalamus (see page 839). Through the fibers of the hippocampal commissure it probably also unites the opposite hippocampal gyri.
Structure of the Cerebral Cortex (Fig. 754).The cerebral cortex differs in thickness and structure in different parts of the hemisphere. It is thinner in the occipital region than in the anterior and posterior central gyri, and it is also much thinner at the bottom of the sulci than on the top of the gyri. Again, the minute structure of the anterior central differs from that of the posterior central gyrus, and areas possessing a specialized type of cortex can be mapped out in the occipital lobe.
On examining a section of the cortex with a lens, it is seen to consist of alternating white and gray layers thus disposed from the surface inward: (1) a thin layer of white substance; (2) a layer of gray substance; (3) a second white layer (outer band of Baillarger or band of Gennari); (4) a second gray layer; (5) a third white layer (inner band of Baillarger); (6) a third gray layer, which rests on the medullary substance of the gyrus.
Nerve Cells.According to Cajal, the nerve cells are arranged in four layers, named from the surface inward as follows: (1) the molecular layer, (2) the layer of small pyramidal cells, (3) the layer of large pyramidal cells, (4) the layer of polymorphous cells.
The Molecular Layer.In this layer the cells are polygonal, triangular, or fusiform in shape. Each polygonal cell gives off some four or five dendrites, while its axon may arise directly from the cell or from one of its dendrites. Each triangular cell gives off two or three dendrites, from one of which the axon arises. The fusiform cells are placed with their long axes parallel to the surface and are mostly bipolar, each pole being prolonged into a dendrite, which runs horizontally for some distance and furnishes ascending branches. Their axons, two or three in number, arise from the dendrites, and, like them, take a horizontal course, giving off numerous ascending collaterals. The distribution of the axons and dendrites of all three sets of cells is limited to the molecular layer.
The Layer of Small and the Layer of Large Pyramidal Cells.The cells in these two layers may be studied together, since, with the exception of the difference in size and the more superficial position of the smaller cells, they resemble each other. The average length of the small cells is from 10 to 15μ; that of the large cells from 20 to 30μ. The body of each cell is pyramidal in shape, its base being directed to the deeper parts and its apex toward the surface. It contains granular pigment, and stains deeply with ordinary reagents. The nucleus is of large size, and round or oval in shape. The base of the cell gives off the axis cylinder, and this runs into the central white substance, giving off collaterals in its course, and is distributed as a projection, commissural, or association fiber. The apical and basal parts of the cell give off dendrites; the apical dendrite is directed toward the surface, and ends in the molecular layer by dividing into numerous branches, all of which may be seen, when prepared by the silver or methylene-blue method, to be studded with projecting bristle-like processes. The largest pyramidal cells are found in the upper part of the anterior central gyrus and in the paracentral lobule; they are often arranged in groups or nests of from three to five, and are named the giant cells of Betz. In the former situation they may exceed 50μ in length and 40μ in breadth, while in the paracentral lobule they may attain a length of 65μ.
Layer of Polymorphous Cells.The cells in this layer, as their name implies, are very irregular in contour; they may be fusiform, oval, triangular, or star-shaped. Their dendrites are directed outward, but do not reach so far as the molecular layer; their axons pass into the subjacent white matter.
There are two other kinds of cells in the cerebral cortex. They are: (a) the cells of Golgi, the axons of which divide immediately after their origins into a large number of branches, which are directed toward the surface of the cortex; (b) the cells of Martinotti, which are chiefly found in the polymorphous layer; their dendrites are short, and may have an ascending or descending course, while their axons pass out into the molecular layer and form an extensive horizontal arborization.
Nerve Fibers.These fill up a large part of the intervals between the cells, and may be medullated or non-medullatedthe latter comprising the axons of the smallest pyramidal cells and the cells of Golgi. In their direction the fibers may be either tangential or radial. The tangential fibers run parallel to the surface of the hemisphere, intersecting the radial fibers at a right angle. They constitute several strata, of which the following are the more important: (1) a stratum of white fibers covering the superficial aspect of the molecular layer (plexus of Exner); (2) the band of Bechterew, in the outer part of the layer of small pyramidal cells; (3) the band of Gennari or external band of Baillarger, running through the layer of large pyramidal cells; (4) the internal band of Baillarger, between the layer of large pyramidal cells and the polymorphous layer; (5) the deep tangential fibers, in the lower part of the polymorphous layer. The tangential fibers consist of (a) the collaterals of the pyramidal and polymorphous cells and of the cells of Martinotti; (b) the branching axons of Golgis cells; (c) the collaterals and terminal arborizations of the projection, commissural, or association fibers. The radial fibers.Some of these, viz., the axons of the pyramidal and polymorphous cells, descend into the central white matter, while others, the terminations of the projection, commissural, or association fibers, ascend to end in the cortex. The axons of the cells of Martinotti are also ascending fibers.
FIG. 754 Cerebral cortex. (Poirier.) To the left, the groups of cells; to the right, the systems of fibers. Quite to the left of the figure a sensory nerve fiber is shown. (See enlarged image)
Special Types of Cerebral Cortex.It has been already pointed out that the minute structure of the cortex differs in different regions of the hemisphere; and A. W. Campbell126 has endeavored to prove, as the result of an exhaustive examination of a series of human and anthropoid brains, that there exists a direct correlation between physiological function and histological structure. The principal regions where the typical structure is departed from will now be referred to.
1. In the calcarine fissure and the gyri bounding it, the internal band of Baillarger is absent, while the band of Gennari is of considerable thickness, and forms a characteristic feature of this region of the cortex. If a section be examined microscopically, an additional layer of cells is seen to be interpolated between the molecular layer and the layer of small pyramidal cells. This extra layer consists of two or three strata of fusiform cells, the long axes of which are at right angles to the surface; each cell gives off two dendrites, external and internal, from the latter of which the axon arises and passes into the white central substance. In the layer of small pyramidal cells, fusiform cells, identical with the above, are seen, as well as ovoid or star-like cells with ascending axons (cells of Martinotti). This is the visual area of the cortex, and it has been shown by J. S. Bolton127 that in old-standing cases of optic atrophy the thickness of Gennaris band is reduced by nearly 50 per cent.
A. W. Campbell says: Histologically, two distinct types of cortex can be made out in the occipital lobe. The first of these coats the walls and bounding convolutions of the calcarine fissure, and is distinguished by the well-known line of Gennari or Vicq dAzyr; the second area forms an investing zone a centimetre or more broad around the first, and is characterized by a remarkable wealth of fibers, as well as by curious pyriform cells of large size richly stocked with chromophilic elementscells which seem to have escaped the observation of Ramón y Cajal, Bolton, and others who have worked at this region. As to the functions of these two regions there is abundant evidence, anatomical, embryological, and pathological, to show that the first or calcarine area is that to which visual sensations primarily pass, and we are gradually obtaining proof to the effect that the second investing area is constituted for the interpretation and further elaboration of these sensations. These areas therefore deserve the names visuo-sensory and visuo-psychic.
2. The anterior central gyrus is characterized by the presence of the giant cells of Betz and by a wealth of nerve fibers immeasurably superior to that of any other part (Campbell), and in these respects differs from the posterior central gyrus. These two gyri, together with the paracentral lobule, were long regarded as constituting the motor areas of the hemisphere; but Sherrington and Grunbaum have shown128 that in the chimpanzee the motor area never extends on to the free face of the posterior central gyrus, but occupies the entire length of the anterior central gyrus, and in most cases the greater part or the whole of its width. It extends into the depth of the central sulcus, occupying the anterior wall, and in some places the floor, and in some extending even into the deeper part of the posterior wall of the sulcus.
3. In the hippocampus the molecular layer is very thick and contains a large number of Golgi cells. It has been divided into three strata: (a) s. convolutum or s. granulosum, containing many tangential fibers; (b) s. lacunosum, presenting numerous vascular spaces; (c) s. radiatum, exhibiting a rich plexus of fibrils. The two layers of pyramidal cells are condensed into one, and the cells are mostly of large size. The axons of the cells in the polymorphous layer may run in an ascending, a descending, or a horizontal direction. Between the polymorphous layer and the ventricular ependyma is the white substance of the alveus.
5. The Olfactory Bulb.In many of the lower animals this contains a cavity which communicates through the olfactory tract with the lateral ventricle. In man the original cavity is filled up by neuroglia and its wall becomes thickened, but much more so on its ventral than on its dorsal aspect. Its dorsal part contains a small amount of gray and white substance, but it is scanty and ill-defined. A section through the ventral part (Fig. 755) shows it to consist of the following layers from without inward:
FIG. 756 Areas of localization on lateral surface of hemisphere. Motor area in red. Area of general sensations in blue. Auditory area in green. Visula area in yellow. The psychic portions are in lighter tints. (See enlarged image)
1. A layer of olfactory nerve fibers, which are the non-medullated axons prolonged from the olfactory cells of the nasal cavity, and reach the bulb by passing through the cribriform plate of the ethmoid bone. At first they cover the bulb, and then penetrate it to end by forming synapses with the dendrites of the mitral cells, presently to be described.
2. Glomerular Layer.This contains numerous spheroidal reticulated enlargements, termed glomeruli, produced by the branching and arborization of the processes of the olfactory nerve fibres with the descending dendrites of the mitral cells.
3. Molecular Layer.This is formed of a matrix of neuroglia, imbedded in which are the mitral cells. These cells are pyramidal in shape, and the basal part of each gives off a thick dendrite which descends into the glomerular layer, where it arborizes as indicated above, and others which interlace with similar dendrites of neighboring mitral cells. The axons pass through the next layer into the white matter of the bulb, and after becoming bent on themselves at a right angle, are continued into the olfactory tract.
4. Nerve Fiber Layer.This lies next the central core of neuroglia, and its fibers consist of the axons or afferent processes of the mitral cells passing to the brain; some efferent fibers are, however, also present, and end in the molecular layer, but nothing is known as to their exact origin.
Weight of the Encephalon.The average weight of the brain, in the adult male, is about 1380 gms.; that of the female, about 1250 gms. In the male, the maximum weight out of 278 cases was 1840 gms. and the minimum weight 964 gms. The maximum weight of the adult female brain, out of 191 cases, was 1585 gms. and the minimum weight 879 gms. The brain increases rapidly during the first four years of life, and reaches its maximum weight by about the twentieth year. As age advances, the brain decreases slowly in weight; in old age the decrease takes place more rapidly, to the extent of about 28 gms.
The human brain is heavier than that of any of the lower animals, except the elephant and whale. The brain of the former weighs from 3.5 to 5.4 kilogm., and that of a whale, in a specimen 19 metres long, weighed rather more than 6.7 kilogm.
Cerebral Localization.Physiological and pathological research have now gone far to prove that a considerable part of the surface of the brain may be mapped out into a series of more or less definite areas, each of which is intimately connected with some well-defined function.
Motor Areas.The motor area occupies the anterior central and frontal gyri and the paracentral lobule. The centers for the lower limb are located on the uppermost part of the anterior central gyrus and its continuation on to the paracentral lobule; those for the trunk are on the upper portion, and those for the upper limb on the middle portion of the anterior central gyrus. The facial centers are situated on the lower part of the anterior central gyrus, those for the tongue, larynx, muscles of mastication, and pharynx on the frontal operculum, while those for the head and neck occupy the posterior end of the middle frontal gyrus.
FIG. 757 Areas of localization on medial surface of hemisphere. Motor area in red. Area of general sensations in blue. Visual area in yellow. Olfactory area in purple. The psychic portions are in lighter tints. (See enlarged image)
Sensory Areas.Tactile and temperature senses are located on the posterior central gyrus, while the sense of form and solidity is on the superior parietal lobule and precuneus. With regard to the special senses, the area for the sense of taste is probably related to the uncus and hippocampal gyrus. The auditory area occupies the middle third of the superior temporal gyrus and the adjacent gyri in the lateral fissure; the visual area, the calcarine fissure and cuneus; the olfactory area, the rhinencephalon. As special centers of much importance may be noted: the emissive center for speech on the left inferior frontal and anterior central gyri (Broca); the auditory receptive center on the transverse and superior temporal gyri, and the visual receptive center on the lingual gyrus and cuneus. |
The Republic of Mozambique is a constitutional democracy with an estimated population of 20 million. President Armando Guebuza was elected in December 2004 in what national and international observers judged to be generally free and fair elections, despite some irregularities. The Front for the Liberation of Mozambique (FRELIMO) has been the ruling political party since independence in 1975, heavily influencing both policymaking and implementation. While civilian authorities generally maintained effective control of the security forces, there were some instances in which elements of the security forces acted independently.
Incidents of serious human rights abuses in some areas—including unlawful killings by security forces and vigilante killings—increased during the year. Prison conditions remained harsh and life threatening, resulting in several deaths. Arbitrary arrest and detention as well as lengthy pretrial detentions were problems. An understaffed and inadequately trained judiciary was inefficient and heavily influenced by the ruling party. Judicial decisions involving independent media outlets created a more constraining environment for press freedom. In addition societal problems such as domestic violence, discrimination against women, abuse, exploitation, and forced labor of children, trafficking in women and children, and discrimination against persons with HIV/AIDS remained widespread.
RESPECT FOR HUMAN RIGHTS
Section 1 Respect for Integrity of the Person, Including Freedom From:
a. Arbitrary or Unlawful Deprivation of Life
Although the government or its agents did not commit any political killings during the year, security forces committed unlawful killings.
Violence as a first resort, excessive use of force, and abuse by police remained problems. Authorities often failed to investigate police violence and bring the perpetrators to justice. However, authorities expelled and, in some cases, brought criminal charges against dozens of officers for disciplinary offenses during the year.
A sharp increase in crime, particularly in and around Maputo City, was likely a factor in the increase in the number of unlawful killings committed by security forces during the year. An overanxious police force responded with a strong show of force and often resorted to violence. Police arbitrarily shot and killed numerous persons. For example, on April 4, three policemen shot and killed Carlos Cossa, Mustafa Assene Momede, and Francisco Antonio Nhantumbo on a football field in a Maputo suburb. An investigation by the Criminal Investigation Police (PIC) indicated the victims were criminals who had escaped from a police car. However, the Attorney General's Office concluded that the shots fired by police were at close range, suggesting summary executions. Further investigation into the killings continued at year's end.
On November 8, members of the riot police shot and killed Juliao Macul in Massinga, Inhambane Province. Minister of Interior Jose Pacheco explained the killing as a police reaction to a denunciation of the suspect and added that Macul attempted to avoid being questioned by police. However, media reports indicated that police mistook Macul for a wanted criminal and made no attempt to identify or arrest Macul before shooting him seven times. Police set up a commission of inquiry to investigate the death, but no further information was available at year's end.
On December 22, police shot and killed Augusto Covilas who had called the police to report that his house was being robbed. Upon arriving at the scene, two members of the police opened fire without attempting to find out who was in the house. Authorities arrested the two officers involved and an investigation was underway at year's end.
Police use of torture resulted in deaths. For example, on August 15, PIC agent Alexandre Francisco Balate drugged, beat, and burned Abranches Afonso Penicelo and left him for dead, according to press reports. Penicelo survived and made it to a hospital before dying from his injuries the following day. Family members of Penicelo claimed the killing was perpetrated by a death squad with which Balate was affiliated. An investigation by the Attorney General's Office was ongoing at year's end.
There were no developments in the following 2006 killings by security forces: the January police killing of an unarmed civilian who was trying to persuade a group of policemen to stop beating a younger woman; the January killing by police of five suspects during a warehouse raid; the March police killing of four detainees in a high security prison; the May killing of at least two prison escapees; the June killing by military police of a secondary student for allegedly wearing boots that belonged to a military officer; the June killing by three members of the Presidential Guard of an unarmed citizen who disobeyed orders to stop his vehicle; and, the July police killings of a prisoner who was active in organized crime and planned to testify in court.
There were no developments in the trial of police who shot six gang members in Matsinho, Manica in 2005.
There was one reported killing as a result of torture and other reports of abuse by members of the Community Policing Councils (CPC), nonstatutory bodies set up by the Mozambican National Police (PRM) in many districts to prevent crime. On April 2, Nzero Zamala, the CPC chairman of Nhamatanda District, Sofala Province, whipped Manuel Chinzo Joao more than 60 times, allegedly because Joao was having an affair with Zamala's wife. Joao was not taken to a hospital until April 6, where he subsequently died from his injuries, according to the Agencia de Informacao de Mocambique (AIM).
Land mine-related accidents resulted in deaths and injuries. The government continued to cooperate with international organizations and donors as well as commercial firms to clear suspected land mine areas.
Unlike in the previous year, there were no reports of killings by unknown actors. The investigation into the March 2006 killing of leading opposition party Mozambican National Resistance (RENAMO) deputy Jose Gaspar Mascarenhas by an unknown gunman was ongoing at year's end.
Killings by vigilante groups were widespread during the year. The nongovernmental organization (NGO) Human Rights League (LDH) and other civil society groups claimed these killings were related to increased crime, lack of police presence in certain neighborhoods, and an ineffective justice system. Most targets of such killings were suspected muggers, thieves, and drug dealers. While nationwide statistics were not available, the press reported at least 26 killings by vigilantes during the year, most of which occurred in and around Maputo City, Matola, and Beira. For example, in June a mob in the Vaz neighborhood of Beira beat to death two suspected thieves. In July a mob killed a suspected thief by beating him, soaking him in gasoline, and burning him to death. In November vigilantes apprehended a suspected petty thief, beat him to death, and burned his body.
Unlike in the previous year, there were no reported mob killings of persons suspected of witchcraft.
There were no reports of politically motivated disappearances.
c. Torture and Other Cruel, Inhuman, or Degrading Treatment or Punishment
While the constitution and law prohibit such practices, police continued to commit abuses. During the year human rights advocates and media outlets reported complaints of torture and other cruel treatment, including several instances involving sexual abuse of women, beatings, and prolonged detention. There were a few reports of deaths resulting from police torture.
The LDH reported that torture in prisons continued to be a problem. Torture and other abusive treatment continued at police squadrons, according to the LDH.
On April 28, police in Macia, Maputo Province, reportedly beat British citizen Alan Evans at a checkpoint. Evans was treated at a clinic in Swaziland for injuries sustained in the incident. There was no further information by year's end.
On July 2, policemen at the 5th squadron in Machava, Maputo Province, severely beat and threatened to kill trainee lawyer Aguinaldo Mandlate, who was at the station to represent clients being interrogated by the PIC. While police claimed he was "trying to escape" and received his injuries after falling, Mandlate stated he fled after police threatened him with a gun. Police gave chase, beat Mandlate, and threw him in a cell with other inmates for several hours before he was rescued by colleagues and taken to the hospital for treatment of head injuries. There were no further updates at year's end.
No action was taken against police involved in the following 2006 torture cases: the May torture of Alexandre Emilio and several others by police from Maputo's 12th Squadron; and the May beating of Generosa Anselmo Cossa, a delegate of RENAMO.
There continued to be reports of abuse and violence by members of the CPC.
Unlike in previous years, there were no reports of violence between FRELIMO and RENAMO supporters during the year.
Vigilante violence also resulted in deaths and injuries.
Prison and Detention Center Conditions
Prison conditions were extremely harsh and life threatening.
The Administration for Prison Services operated 211 prisons in 10 provinces and is under the Ministry of Justice. The Ministry of Interior and the police are responsible for jails at police stations.
Overcrowding remained a serious problem. The LDH noted that many prisons held up to three times the number of prisoners for which they were built, and that often prisoners slept in bathrooms, standing up, or in shifts. For example, the Maputo Central Prison held 2,246 prisoners in a facility designed to hold 800, and the Inhambane Provincial Prison held 346 prisoners in a facility designed to hold 75. During the first half of the year, the LDH visited 74 prisons and detention facilities, which held a total of 11,424 inmates in facilities designed to hold 5,913.
The LDH found that more than 500 detainees in the Maputo Central Prison (Machava) had been held beyond the 90-day preventive detention period. Of the prisons visited, 399 prisoners remained in jail after the end of their sentences (including 206 at the Maputo Central Prison). The LDH described 35 facilities as "physically inadequate." In detention facilities, overcrowding did not appear to be a serious problem. During the first half of the year, the LDH visited several police station detention facilities and noted that some detainees continued to be held beyond the maximum police station preventive detention period of 48 hours.
Reports continued that most prisoners received only one meal per day. In 13 of the prisons visited, the LDH characterized the provision of food as "poor." It was customary for families to bring food to prisoners; however, there continued to be occasional reports that guards demanded bribes in exchange for delivering food to prisoners. In several prisons, inmates prostituted themselves in exchange for food, according to the LDH.
There continued to be many reported deaths in prison, the vast majority due to illness. In many facilities overcrowding, lack of sanitation, potable water, and food also led to sickness.
In a series of prison visits conducted during the first half of the year, the LDH found malaria, scabies, and tuberculosis to be frequent among prisoners in nearly all of the country's prisons. LDH also found other illnesses caused by malnutrition, including paralysis and blindness. Both healthy and sick prisoners regularly were kept in the same cells. The spread of HIV/AIDS and other sexually transmitted diseases was a serious problem for the prison population, and the LDH noted that in many prisons, authorities denied condoms to inmates.
In the first half of 2007, the LDH reported 39 juveniles under the age of 16 held with adults from the general prison population.
Pretrial detainees were held with convicted prisoners.
International and domestic human rights groups had access to prisoners, although at the discretion of ministries of justice and interior. Unlike in previous years, the LDH reported facing serious problems obtaining credentials to visit prisons. While in some cases authorities simply did not respond to the LDH requests for credentials, in other cases certain prison directors denied the LDH access to visits despite having credentials.
d. Arbitrary Arrest or Detention
While the constitution and law prohibit arbitrary arrest and detention, both practices continued to occur.
Role of the Police and Security Apparatus
Forces under the Ministry of Interior, including the PIC, the PRM, and the Rapid Intervention Force, are responsible for internal security. An additional security body, the State Information and Security Service, reports directly to the president. The armed forces (FADM) are responsible for external security, but patrolled with the PRM and manned checkpoints due to a significant increase in crime during the year.
The police continued to be poorly paid, despite an increase in pay during the year. Trainee-level officers reportedly received approximately $89 (2,113 meticais) a month, while those at higher rank received approximately $115 (2,725 meticais) a month. Corruption and extortion by police were widespread, and impunity remained a problem. In January Deputy Interior Minister Jose Mandra said that criminal elements had infiltrated the police force.
Police regularly detained persons for arbitrary reasons and demanded identification documents solely to extort payments. Many crime victims reportedly avoided police assistance because of expected demands for bribes and a lack of confidence that the police would help. During the 12 months preceding April, the Maputo City Police Command initiated disciplinary and criminal proceedings against 113 Maputo policemen, expelling 28 of these from the force. The most common reasons for disciplinary action, according to Maputo's police chief, were collaboration with criminals, extortion of goods and money, excessive alcohol consumption, and abandonment of post. During the year the Ministry of Interior expelled at least 160 police officers. However, the vast majority of police who committed infractions were "recycled," sent back to school, and then transferred to a new unit. In the three months preceding March, the Ministry "recycled" 178 police. These included suspected criminals, thieves, and agents suspected of collaborating with criminals. A government-sponsored survey ranked the PRM as the second most corrupt public institution.
Professional training for police officers continued during the year; in August 60 PRM officers in Gaza Province completed human rights training.
Implementation of the 2003-12 strategic plan of action and modernization for the PRM continued; seven of its nine "guiding principles" reflected respect for human rights. While the plan acknowledged the problem of abuse of police powers, it made no specific provision for ensuring greater accountability for such abuses.
Arrest and Detention
Although the law provides that persons be arrested openly with warrants issued by a judge or prosecutor (except persons caught in the act of committing a crime), police continued to arrest and detain citizens arbitrarily. By law the maximum length of investigative detention without a warrant is 48 hours, during which time a detainee has the right to judicial review of the case. The individual may be detained another 90 days while the PIC continues its investigation. When a person is accused of a crime carrying a sentence of more than eight years, the individual may be detained up to 84 days without being charged formally. With court approval, such detainees may be held for two more periods of 84 days each without charge while the police complete the investigative process. The law provides that when the prescribed period for investigation has been completed and no charges have been brought, the detainee must be released. In many cases the authorities either were unaware of these regulations or ignored them, often also ignoring a detainee's constitutional right to counsel and to contact relatives or friends. The law provides that citizens have access to the courts as well as the right to representation, regardless of ability to pay for such services. However, due to a shortage of legal professionals, indigent defendants frequently had no legal representation.
The bail system remained poorly defined. Prisoners, their families, and NGOs continued to complain that police and prison officials demanded bribes for releasing prisoners.
There were reports that police harassed and arbitrarily detained persons, including journalists, during the year.
Government statistics indicated that approximately 40 percent of inmates were still awaiting trial. In addition there continued to be reports of detainees who spent longer in pretrial detention than the period of the sentence they eventually received. By law a judge has 48 hours to validate a detention in any proceeding; however, this statute often was not enforced.
e. Denial of Fair Public Trial
The constitution and law provide for an independent judiciary; however, the executive branch and the ruling FRELIMO party heavily influenced an understaffed and inadequately trained judiciary, particularly in the lower tiers. The judicial system continued to suffer from a lack of transparency and often did not comply with the principles of promotion and protection of human rights.
In May the parliament passed a new judicial organization law, which establishes intermediate appeals courts and expands the powers of district courts to rule on more serious criminal cases. The new law also empowers district court judges to rule on criminal cases with penalties ranging between eight and 12 years, compared with up to only two years before the law. The law permits judges to rule on a significantly higher number of cases and was expected to reduce the backlog in the judicial sector. In addition alternative measures such as work brigades, conditional release for prisoners who have completed half of their sentence, and traveling tribunals continued.
Approximately 93 of the country's 128 judicial districts had functioning courts; however, a shortage of judges and qualified staff was a major problem. In March Chief Justice Mario Mangaze reiterated that the country had only 36 percent of the judges and prosecutors needed to administer justice effectively. There were 221 judges (or approximately one per 90,500 inhabitants), 183 of whom held law degrees as required by law for all judges appointed after 2000. During the year 7 percent of the 1,429 staff employed by the courts held university degrees. Continuing problems included chronic absenteeism, unequal treatment, low salaries, corruption, deliberate delays, and omissions in handling cases. Unlike in the previous year, there were no reports that judges were expelled for illicit behavior.
The president appoints both the Supreme Court president and vice president. The Higher Judicial Magistrates' Council (CSMJ) prepares Supreme Court nominations and submits a list of qualified potential nominees to the president. Members of the CSMJ tended to be either FRELIMO members or FRELIMO-affiliated. The president makes all other judicial appointments.
There are two complementary formal justice systems: the civil justice system and the military justice system. The Supreme Court administers the civil system, and the Ministry of National Defense administers military courts. Under the Supreme Court there are province and district-level courts, and each province has a court of appeal. Cases in military courts may be appealed to the Supreme Court. Civilians are not under the jurisdiction of, or tried in, the military courts.
There also are courts that exercise limited, specialized jurisdiction, such as the Administrative Court, the Customs Court, and the Maritime Court. The Constitutional Council is charged with determining the constitutionality of laws and decrees, supervising the electoral process, declaring and validating electoral results, and ruling on electoral disputes. A separate court system exists for minors 16 years of age and younger. The government may send minors to correctional, educational, or other institutions.
Persons accused of crimes against the government are tried publicly in regular civilian courts under standard criminal judicial procedures. Members of the media may attend trials, although space limitations prevented the general public from attending. A judge may order a trial closed to the media in the interest of national security or to protect the privacy of the plaintiff in a sexual assault case. Article 12 of the judicial organization law "prohibits the production and public transmission of images and sounds at trials." There is no trial by jury.
In regular courts all accused persons, in principle, are presumed innocent and have the right to legal counsel and appeal; however, authorities did not always respect these rights. Although the law specifically provides for public defenders for the accused, such assistance generally was not available in practice, particularly in rural areas. The LDH reported that most citizens remained unaware of this right, and many had no access to legal counsel. Some NGOs continued to offer limited legal counsel at little or no cost to both defendants and prisoners. Only judges or lawyers may confront or question witnesses.
Outside the formal court system, local customary courts, and traditional authority figures often adjudicated matters such as estate and divorce cases. Respected local arbiters with no formal training staffed customary courts.
Political Prisoners and Detainees
There were no confirmed reports of political prisoners or detainees. Unlike in the previous year, RENAMO did not continue to allege that 10 of its party members were being held as political prisoners in Mutarara District in Tete Province.
Civil Judicial Procedures and Remedies
Although the law provides for an independent and impartial judiciary in civil matters, in practice the judiciary was subject to political interference.
f. Arbitrary Interference with Privacy, Family, Home, or Correspondence
The constitution and law prohibit such actions, and the government generally respected these prohibitions in practice. However, opposition party members alleged that government intelligence services and ruling party activists continued without warrants to monitor telephone calls, conduct surveillance of their offices, follow opposition members, use informants, and disrupt party activities in certain areas of the country, including in Cabo Delgado and Nampula Provinces. By law police require a warrant to enter homes and businesses and also to monitor telephone calls.
In August the media reported that the FADM required senior members to complete a confidential questionnaire regarding party affiliation, activities in support of the party, and whether the individual supported FRELIMO. Some observers believed the questionnaire was evidence of the politicization of the FADM. Media reports noted that while Minister of Defense Tobias Dai denied knowledge of the questionnaire, Brigadier General Jorge Gune acknowledged that he had filled out several questionnaires since achieving higher rank in the FADM.
Unlike in the previous year, there were no reports that school administrators in Muecate district in Nampula Province forced single male teachers to marry to mitigate the number of sexual assaults of female students by teachers.
Section 2 Respect for Civil Liberties, Including:
a. Freedom of Speech and Press
Although the constitution and law provide for freedom of speech and of the press, in practice there were some restrictions on these rights. There were occasional reports that police harassed journalists, and journalists admitted that self-censorship was common. In its annual report on press freedom, the NGO Freedom House noted that the press was "partially free," while also acknowledging the continued growth of private media. The NGO Media Institute of Southern Africa (MISA) noted in its annual report that court decisions involving several independent media outlets during the year created a constraining environment. Individuals could criticize the government publicly or privately without reprisal.
The independent media were active and expressed a wide variety of views. The government maintained majority ownership in Noticias, the main newspaper and the only daily with nationwide distribution. Noticias, the daily Diario de Mocambique, and the weekly Domingo largely reflected the views of the government and provided marginal, often critical coverage of RENAMO but also demonstrated a willingness to examine government actions. The government-run news agency, AIM, often printed stories critical of the performance of government ministries or agencies.
The international media were allowed to operate freely.
There were numerous private radio stations that operated throughout the country. Radio Mocambique, which received 60 percent of its operating budget from the government, was the most influential media service with the largest audience in the country. While broadcasting debates on important issues of the country, Radio Mocambique tended to invite participants that were not critical of the government.
MISA noted that the process for obtaining a radio operating license was often long, convoluted, and politically biased. According to MISA, a new law was needed which would clearly delineate the difference between commercial and public radio.
The government supplied 80 percent of the operating budget for Televisao de Mocambique (TVM), the television station that broadcast to the largest percentage of the population. While TVM provided more balanced news coverage than in previous years, it retained a strong government and FRELIMO bias.
On January 3, the Maputo City Court ordered the return of all equipment seized by law enforcement officers in December 2006 from the private television station STV. While the seizure ostensibly involved a severance pay dispute, STV was a frequent critic of the government, leading many civil society groups to believe the seizure was a crackdown on STV.
Security forces harassed and arbitrarily detained local journalists during the year. In January police in Beira detained without charge photojournalist Celeste MacAurthur of the Diario de Mocambique for taking pictures of an abandoned house. Police released MacAurthur later the same day. In March police arrested Celso Manguana, a reporter for Canal de Mocambique, and detained him for three days. Police accused Manguana of "insulting authority" after he went to a police station to inquire about the arrests of several demonstrators. He was released after the intervention of the attorney general and the LDH.
In August Canal de Mocambique reported that one of its journalists, Luis Nhachote, received a threat message on his cell phone apparently for an article he published critical of FRELIMO.
In November Nyimpine Chissano, oldest son of former President Joaquim Chissano, and "joint moral author" in the killing of investigative journalist Carlos Cardoso in 2000, died.
The National Union of Journalists, MISA, and the Mozambican Editor's Forum criticized the new judicial organization law. The groups argued that Article 12, which "prohibits the production and public transmission of images and sounds at trials" was a serious threat to press freedom by imposing a blanket ban on microphones and cameras in the courtroom.
Defamation of the president is prohibited; however, no one was charged with the offense during the year.
MISA condemned at least two court rulings in libel suits brought against the independent newspapers Horizonte and Faisca during the year. While the cases did not involve the government, MISA argued that the damages demanded would shut down the newspapers and silence the few alternatives to government sources of information.
Newsprint and other printing supplies must be imported from South Africa, and the government did not exempt these supplies from import duties. Some newspapers found it more cost-effective to print in South Africa and import the final product. Other journals were only published in electronic versions, severely limiting their readership. Journals printed on paper had limited readership beyond Maputo, due to high transportation costs.
There were no government restrictions on access to the Internet or reports that the government monitored e-mail or Internet chat rooms. Individuals and groups could engage in the peaceful expression of views via the Internet, including by e-mail. While public access to the Internet continued to expand, particularly in the larger cities, lack of infrastructure in the rural parts of the country and installation costs limited overall use.
Academic Freedom and Cultural Events
While the government generally did not restrict academic freedom, there were reports that teachers at the university, secondary, and primary school level felt pressure to align themselves with FRELIMO, particularly in the central and northern provinces.
In April the Administrative Tribunal reinstated Ismael Mussa, a RENAMO parliamentary deputy and lecturer demoted from his position as director of social services at the state-run Eduardo Mondlane University (UEM) in 2005. Despite the ruling, in July the UEM again removed Mussa from his position as director. While university regulations allow the vice chancellor to appoint or dismiss directors, observers suspected political harassment.
b. Freedom of Peaceful Assembly and Association
Freedom of Assembly
The constitution and law provide for freedom of assembly; however, there was one instance in which police briefly detained demonstrators during the year. While the law regulates public demonstrations, it does not apply to private gatherings held indoors and by individual invitation, nor does it affect religious gatherings or election campaigning.
In April police, apparently without cause, detained 10 protestors in front of the National Assembly building. According to the newsfax A Canal de Mocambique, six of the protestors were held overnight before being released without explanation the following day.
In March, for the first time since 2004, local authorities in Maputo city permitted a group of madjermanes to hold a public march. The madjermanes, a group of approximately 15,000 citizens who worked in the former East Germany, demanded payment of benefits for their past work.
Freedom of Association
The constitution and law generally provide for freedom of association, although the government imposed some limits on this right. According to the law a political party is required to demonstrate that it has no regional, racial, ethnic, or religious exclusiveness and must secure at least 2,000 signatures to be recognized. There were approximately 50 registered political parties.
A government decree regulates the registration and activities of foreign NGOs. The registration process for foreign NGOs and religious groups reportedly involved significant discretion on the part of government officials and regularly took several months.
c. Freedom of Religion
The constitution and law provide for freedom of religion, and the government generally respected this right in practice.
The constitution and the law governing political parties specifically forbid religious groups from organizing political parties and any political party from sponsoring religious propaganda as threats to national unity.
The Catholic Church and some Muslim communities continued to request the return of certain properties nationalized by the government in the years immediately following independence, including schools, health centers, shops, and residences.
Societal Abuses and Discrimination
Relations among various religions groups were generally amicable. In September the PRM arrested Fernando Bernardo Arrone and Fabiao Domingos for their roles in burning three mosques in late August in Lichinga, Niassa Province. Arrone confessed that he received $337 (8,000 meticais) from the Catholic Church for each mosque. The investigation into the incidents was ongoing at year's end.
There was a very small Jewish population, and there were no reports of anti-Semitic acts.
For a more detailed discussion, see the 2007 International Religious Freedom Report.
d. Freedom of Movement, Internally Displaced Persons, Protection of Refugees, and Stateless Persons
While the law provides for freedom of movement within the country, foreign travel, emigration, and repatriation, the government sometimes infringed upon these rights in practice.
Traffic checkpoints are legal and under the jurisdiction of traffic police. Checkpoints occasionally affected freedom of movement, and according to press reports, authorities sometimes abused and demanded bribes from citizens at checkpoints. Police sometimes stopped foreigners and ordered them to present original passports or resident papers, refused to accept notarized copies, and fined or detained those who failed to show proper documents. Police, including members of CPCs, also routinely harassed, detained, and extorted bribes from local citizens for failure to carry identity papers.
The law prohibits forced exile, and the government did not use it.
Protection of Refugees
The laws provide for the granting of asylum or refugee status in accordance with the 1951 UN Convention relating to the Status of Refugees and its 1967 protocol, and the government has established a system for providing protection to refugees. In practice the government provided protection against "refoulement," the return of persons to a country where there is reason to believe they feared persecution. The government cooperated with the UN High Commissioner for Refugees (UNHCR) and other humanitarian organizations in assisting refugees and asylum seekers.
While the government assisted in the repatriation of 300 refugees in May, Canal de Mocambique reported in July that due to a lack of resources, the government was not able to satisfy requests by some refugees to return to their countries of origin.
In April AIM reported on several attacks against Burundian and Congolese refugees in Nampula Province. The government continued to limit refugee movement within the country. Refugees must request authorization to move outside the geographic region in which they have been registered. In addition refugees residing within the Marratane camp in Nampula Province must request authorization to leave its boundaries, which has perpetuated the extracting of bribes by officials.
Section 3 Respect for Political Rights: The Right of Citizens to Change Their Government
The constitution and law provide citizens the right to change their government peacefully, and citizens exercised this right in practice through periodic, free, and fair elections held on the basis of universal suffrage.
Elections and Political Participation
In 2004 citizens elected Armando Guebuza of the ruling FRELIMO party as president in the country's third multiparty general elections. While domestic and international observers noted that voting day procedures generally followed international norms, they also documented irregularities during the campaign and in the vote count. FRELIMO used significant state funds and resources for campaign purposes, in violation of election law. RENAMO issued complaints of election fraud to several agencies, including the Constitutional Council. In January 2005 the Constitutional Council affirmed Guebuza as the winner.
Unlike in previous years, there were no reports of violence between FRELIMO and RENAMO supporters during the year.
There were 93 women in the 250-seat National Assembly. The prime minister was a woman, and women held six of the 24 ministerial positions and four of the 18 vice ministerial positions. Women held 30 percent of the seats on FRELIMO's 160-member Central Committee and six seats on the 17-member Political Commission.
Members of many ethnic groups held key positions in both the legislative and executive branches. There was no evidence that specific ethnic groups were excluded.
Government Corruption and Transparency
While the law provides criminal penalties for official corruption, the government did not implement the law effectively, and officials often engaged in corrupt practices with impunity. No corruption cases involving high-profile individuals have been brought to trial during the Guebuza administration.
The World Bank's Worldwide Governance Indicators reflected that corruption was a serious problem.
Despite the government's strong anticorruption rhetoric, corruption in the executive and legislative branches was widely perceived to be endemic. Petty corruption by low-level government officials to supplement low incomes, and high-level corruption by a small group of politically connected elites continued to be the norm. Corruption largely resulted from a lack of checks and balances, minimal accountability, and a culture of impunity. Local NGOs, such as the Center for Public Integrity, and media groups continued to be the main forces fighting corruption, reporting and investigating numerous corruption cases. The law requires that all members of the government declare and deposit their assets with the Constitutional Council, but does not require that such information be made available to the general public.
The Central Office for the Combat of Corruption (GCCC) functions as an autonomous unit under the attorney general's office with its own state budget. According to the GCCC, from January to August prosecutors brought charges in 13 cases of corruption. In December the Ministry of Civil Service reported that authorities expelled nearly 400 public servants for various irregularities during the year. In August the Supreme Court refused to consider some 15 corruption cases brought forward by the GCCC after several judges claimed the GCCC lacked legal authority to prosecute. In December the attorney general announced that the GCCC would no longer have the power to investigate cases of forgery, swindling, murder, and theft, drastically reducing its scope.
Several new cases of corruption were reported. In January the GCCC ordered the arrest of Deputy Director of the Maputo Central Prison Arminda Parruque for under invoicing and the disappearance of large sums of money. Parruque was being held at the Maputo Civil Prison, and an investigation was ongoing at year's end.
In January authorities arrested six health service administrators in Cabo Delgado Province for the theft of approximately $126,000 (3 million meticais) intended for funding health services in six districts. There were no further developments by year's end.
In March authorities arrested the provincial administrator for youth and sports in Niassa Province for the theft of $76,000 (1.8 million meticais). The investigation was ongoing at year's end.
There were no further developments in the 2006 investigations into alleged corruption by government officials.
The NGO Etica Mozambique, which operated corruption reporting centers in major cities to provide a mechanism for citizens to anonymously report incidences of corruption, became inactive during the year. Management and resource constraints were ongoing problems, and none of the cases it had transferred to the Ministry of Justice had gone to trial.
There were no laws providing for public access to government information, and in practice the government restricted citizens' and noncitizens' access to public information.
Section 4 Governmental Attitude Regarding International and Nongovernmental Investigation of Alleged Violations of Human Rights
Domestic and international human rights groups generally operated without government restriction, investigating and publishing their findings on human rights cases. Although at times slow, government officials were generally cooperative and responsive to their views. Registration procedures for NGOs often were lengthy.
While an independent ombudsman position to investigate allegations of abuses, including human rights violations, by state officials was created by constitutional amendment in 2005, an ombudsman had yet to be named.
Section 5 Discrimination, Societal Abuses, and Trafficking in Persons
The constitution and law prohibit discrimination based on race, gender, disability, language, or social status, but in practice discrimination persisted against women, persons with disabilities, and persons with HIV/AIDS.
The law prohibits rape (excluding spousal rape) but was not effectively enforced. Penalties ranged from two to eight years' imprisonment if the victim is 12 years of age or older, and eight to 12 years' imprisonment if the victim is under the age of 12. While there were no official estimates as to the extent of spousal rape, it was regarded as a common problem. According to NGO reports, many families preferred to settle such matters privately through financial remuneration rather than through the formal judicial system. In August the Maputo High Court sentenced Luis Camillo to 17 years' imprisonment for the 2004 rape of two South African women near Johannesburg.
Reports indicated that domestic violence against women, particularly spousal rape and beatings, was widespread, and the PRM received 5,667 reports of violence against women through September. There is no law that defines domestic violence as a crime, but laws prohibiting rape, battery, and assault could be used to prosecute domestic violence. In many circles women believed it was acceptable for their husbands to beat them. Cultural pressures discouraged women from taking legal action against abusive spouses.
A 15-month survey released in August 2006 revealed that 54 percent of women respondents admitted suffering an act of physical or sexual violence by a man at some point in their lives, 37 percent in the last five years, and 21 percent during the past year.
There was no update on the December 2006 case of Antineco Chibewa, who killed his 36-year-old wife for being too old.
The government and NGOs often worked together to combat domestic violence. The PRM operated special women and children's units in police squadrons that received cases of domestic violence, sexual assault, and violence against children; the units provided assistance to victims and their families. All 30 police squadrons in Maputo had women and children's centers. In addition all police squadrons in the country installed a "green line" (a free phone line) to receive complaints of violence against women and children.
Kukuyana, a national network of women living with HIV/AIDS, reported that many women were expelled from their homes and/or abandoned by their husbands and relatives because they were HIV positive. It also reported that some women who were widowed by HIV/AIDS were accused of being witches who purposely killed their husbands to acquire belongings, and in retribution were deprived of all belongings.
Prostitution is legal, although it is governed by several laws against indecency and immoral behavior and restricted to certain areas. The practice was widespread and particularly prevalent along major transportation corridors and in border towns where long-distance truckers overnight. Young women without means of support were at the greatest risk for being drawn into prostitution.
Sexual harassment is illegal; however, it was pervasive in business, government, and education. Although no formal data exists, the media reported numerous instances of harassment during the year.
Forced marriage of girls and women was a problem.
"Purification," whereby a widowed woman is obligated to have unprotected sex with a member of her husband's family, continued to be practiced, particularly in rural areas.
With the exception of some ethnic and religious groups, the groom's family provided a dowry to the bride's family, usually in the form of livestock, money, or other goods. For Muslims, the bride's family usually paid for the wedding and provided gifts. These exchanges contributed to violence and other inequalities, due to the perception that the women subsequently were "owned" by the husband.
The Family Law (which took effect in 2005) sets the age of marriage for both genders at 18 for those with parental consent, and 21 for those without parental consent. The law also eliminates husbands' de facto status as heads of families, and legalizes civil, religious, and common law unions. While the law does not recognize new cases of polygyny, it grants women already in polygynous marriages full marital and inheritance rights. The law more precisely defines women's legal rights with regard to property, child custody, and other issues. However, nearly three years after taking effect, a survey conducted by the NGO MULEIDE found that approximately 63 percent of women remained uninformed about the law. A Save the Children report on inheritance practices released in June noted that 60 percent of women cited discrimination in the inheritance process. The same report highlighted cases in which women lost inheritance rights for not being "purified" following the death of their husbands.
Customary law was still practiced in many parts of the country. In some regions, particularly the northern provinces, women had limited access to the formal judicial system for enforcement of rights provided under the civil code and instead relied on customary law to settle disputes. Under customary law, women have no rights to the disposition of land.
The law grants citizenship to the foreign-born wife of a male citizen but not to the foreign-born husband of a female citizen.
Women continued to experience economic discrimination, were three times less likely than men to be represented in the public and private sectors, and often received lower pay than men for the same work.
While the government continued to stress the importance of children's rights and welfare, significant problems remained.
A UN Children's Fund (UNICEF) report released in May estimated that the level of birth registration was less than 40 percent, and that 94 percent of children under age four were not registered. In some cases, particularly in rural areas, lack of birth certificates prevented citizens from registering for school, access to health care, and the right to obtain public documents, such as identity cards or passports.
Education is compulsory through age 12, but enforcement was inconsistent due to the lack of resources and the need for additional schools. Public education is free, but most families paid enrollment fees for each child and purchased books, uniforms, and other school supplies. Children who have a certificate that testifies their parents' incomes are below a certain poverty level are exempt from fees, but for most families, fees and associated costs remained a significant financial burden.
During the year UNICEF estimated that 94 percent of children were enrolled in school; primary school enrollment reached 4.5 million, and secondary school enrollment increased from 45,000 to approximately 360,000 since 1992. Despite joint government/NGO initiatives in specific localities and districts to improve girls' school attendance, completion rates for primary school students were approximately 41 percent for boys and 29 percent for girls. In January a report released by Save the Children noted that more than one million children (the majority of whom were girls) between the ages of six and 11 either had never been to school or did not currently attend. Primary schools remained overcrowded, and approximately 70 percent lacked adequate sanitation.
The PRM reported more than 2,800 cases of child abuse during the year, but noted that the vast majority of cases went unreported. Most cases involved sexual abuse, physical abuse, or negligence. Several cases of fathers sexually abusing their daughters were reported during the year. Sexual abuse in schools was a growing problem. An analysis undertaken by Actionaid International in Zambezia, Manica, and Maputo provinces revealed that 78 percent of girls between eight and 18 were forced to have sex with their professors to pass their class. A separate joint study by UNICEF, Actionaid International, and Save the Children revealed that one in five girls over the age of 15 reported being sexually abused by professors and that most children did not report these cases because they were afraid or ashamed. The press continued to report cases in which primary and secondary school students often paid teachers in exchange for a spot in a class or better grades.
There continued to be reports in newspapers of physical abuse of students by teachers during the year.
Local customs, primarily in the northern provinces and in Muslim and South Asian communities, resulted in underage marriage. The daily Noticias reported that in the rural areas of Nampula Province, some districts reported 10 percent fewer female students enrolled in school compared with 2006 as a result of child marriage.
While the law prohibits pornography, child prostitution, and sexual abuse of children under 16 and proscribes prison sentences and fines for perpetrators, exploitation of children below the age of 15 continued, and child prostitution remained a problem. In practice perpetrators of these crimes rarely were identified and prosecuted, and punishment was not commensurate with the crime.
The country continued to have a problem with street children, but no nationwide figures were available.
Zimbabwean children, many of whom entered the country alone, continued to face labor exploitation and discrimination. They lacked protection due to inadequate documentation and had limited access to schools and other social welfare institutions. Coercion of girls into the sex industry was common.
The government took steps to address the problems facing HIV/AIDS orphans. A 2006 UNICEF study estimated that of the country's 1.6 million orphans, more than 380,000 lost either one or both parents to HIV/AIDS. Several government agencies, including the Ministry of Health and the Ministry of Women and Social Action, implemented programs to provide health assistance and vocational education for HIV/AIDS orphans.
The Maputo City Office of Women and Social Action continued its program of rescuing abandoned orphans and assisting single mothers who headed families of three or more persons. They also offered special classes to children of broken homes in local schools. NGO groups sponsored food, shelter, and education programs in all major cities.
Trafficking in Persons
Although the law does not prohibit trafficking in persons, traffickers could be prosecuted using 13 related articles of the penal code on sexual assault, rape, abduction, and child abuse.
The country was a source, transit, and possibly a destination for trafficked persons. While there were no official statistics, NGOs believed that trafficking was becoming a serious problem. UNICEF released a report in January that revealed more than 1,000 cases of women and children trafficked from Mozambique to South Africa between 2002 and 2006. Most trafficking victims were transported to South Africa on the highway from Maputo to Johannesburg. The majority of victims were women and children trafficked for both sexual exploitation and forced labor. Boys were trafficked as laborers on South African farms and in mines, and girls were trafficked to work as prostitutes and as domestics. Poverty, a history of child migration, and weak border controls all contributed to trafficking.
Child prostitution appeared to be most prevalent in Maputo, Nampula, Beira, and at border towns and overnight stopping points along key transportation routes. Child prostitution reportedly was growing in the Maputo, Beira, and Nacala areas, which had highly mobile populations and a large number of transport workers. Child prostitution also was reported in Sofala and Zambezia provinces. Some NGOs provided health care, counseling, and training in other vocations to children engaged in prostitution.
Traffickers were principally citizens or South African. Trafficking groups included small networks of citizens based in Maputo and Nampula, and there were reports that organized crime groups were involved. Traffickers often lured victims by promising better jobs in South Africa. Once there, they were threatened with exposure of their illegal status and forced to work for little or no pay. Often women were sexually assaulted en route to their destination or once they arrived in South Africa. There were also reports that syndicates trafficked young girls from Thailand through the country en route to South Africa.
The government's law enforcement efforts decreased over the previous year, and a paucity of training resources continued to hinder greater efforts. There were no prosecutions or convictions for trafficking cases during the year. Many lower-ranking police and border control agents were suspected of accepting bribes from traffickers.
Due to a lack of resources, government officials regularly called on NGOs for the provision of protection and assistance to victims, including shelter, food, counseling, and rehabilitation. The Ministry of Interior expanded the number of offices for attending to women and child victims of violence from 96 to 152, and provided victims' assistance training for police officers who dealt with such cases. The police also conducted general training on trafficking and detecting at-risk children in the central provinces of Sofala, Manica, and Zambezia and the northern province of Nampula.
Persons with Disabilities
Although the constitution and law stipulate that citizens with disabilities shall fully enjoy the same rights as all other citizens, the government provided few resources to implement this provision. Discrimination was common against persons with disabilities in employment, education, access to health care, and in the provision of other state services. The law does not mandate access to buildings for persons with disabilities, but the Ministry of Public Works and Habitation worked to ensure that public buildings in Maputo city provided access to persons with disabilities. Electoral law provides for the needs of voters with disabilities in the polling booths.
Concerns of persons with disabilities included lack of access to socioeconomic opportunities and employment, limited access to buildings and transportation, and a lack of wheelchairs. Special access facilities were rare. There were few job opportunities for persons with disabilities in the formal sector.
The country's only psychiatric hospital was overwhelmed with patients and lacked the means to guarantee even basic nutrition, medicine, or shelter. During the first six months of the year, the hospital received 1,160 patients, compared with 348 during the same period in 2006. Doctors at the hospital also reported that many abandoned family members with disabilities at the hospital. Veterans with disabilities continued to complain about not receiving pensions.
The Ministry of Women and Social Action is responsible for protecting the rights of persons with disabilities. The four-year National Action Plan in the Area of Disabilities announced in 2006 still required budget allocation to be effectively implemented.
Maputo city offered free bus passes to persons with disabilities.
There were reports of tension between newly arrived Chinese guest workers, often employed in construction, and citizens in Maputo city and Beira, Sofala Province.
There were reports of discrimination by police against Zimbabwean immigrants during the year.
Unlike in the previous year, there were no reports of vigilante killings of West African immigrants.
Other Societal Abuses and Discrimination
The law prohibits discrimination on the basis of HIV/AIDS, and the Ministry of Labor generally intervened in cases of perceived discrimination by employers. In July the Ministry of Labor reported receiving more than 100 cases annually of workers being dismissed by their employers for having HIV/AIDS. Often, the worker was obligated by the employer to take HIV/AIDS tests. In response to these violations, the ministry registered the complaints and confronted companies responsible for dismissals.
The law does not specifically prohibit discrimination based on sexual orientation, and there were occasional such reports. Despite the absence of a law, the LDH reported cases of discrimination against homosexuals in the judicial system. The Workers Law, passed during the year, includes an article that prevents discrimination in the workplace based on a number of factors, including sexual orientation.
Section 6 Worker Rights
a. The Right of Association
The constitution and law provide that all workers are free to join a trade union of their choice without previous authorization or excessive requirements, and workers exercised these rights in practice. Labor laws guaranteeing the right of association do not cover government employees, including firefighters, members of the judicial authorities, and prison guards. As of June the Mozambican Workers' Association (OTM) estimated that of the approximately 500,000 workers in the formal sector, 98,000 were unionized. Some unions alleged that the OTM was under the influence of FRELIMO.
The law prohibits antiunion discrimination; however, there were reports that many companies continued to engage in antiunion discrimination by replacing persons at the end of contracts, dismissing workers for going on strike, and not abiding by collective bargaining agreements.
b. The Right to Organize and Bargain Collectively
Although the law provides for the right of workers to organize and engage in collective bargaining, such contracts covered less than 2 percent of the work force. The government did not set private sector salaries. Unions were responsible for negotiating wage increases.
The law explicitly provides for the right to strike, and workers exercised this right in practice; however, civil servants, police, military personnel, and workers in other essential services (including sanitation, firefighting, and health care) do not have the right to strike. The law specifies that strikers must notify police, the government, union, and employers 48 hours in advance of intended strikes.
On July 16, the head of the Mafambisse security force in Sofala province shot and killed striking worker Domingos Chanjane and injured two others. While workers participating in the strike insisted the perpetrator was also a member of the police, a PRM spokesman denied the claim. There were no further updates at year's end.
There are no special laws or exemption from regular labor laws in the few export processing zones.
c. Prohibition of Forced or Compulsory Labor
The law prohibits forced and compulsory labor, including by children, and while there were few reports that such practices occurred in the formal economy, forced and bonded labor, particularly by children, was common in rural areas.
d. Prohibition of Child Labor and Minimum Age for Employment
While the law prohibits child labor, it remained a problem. In the formal economy, the minimum working age without restrictions is 18 years of age. The law permits children between 15 and 18 to work, but the employer is required to provide for their education and professional training and to ensure that conditions of work are not damaging to their physical and moral development. Children between the ages of 12 and 15 are permitted to work under special conditions authorized jointly by the ministries of labor, health, and education. For children under the age of 18, the maximum workweek is 38 hours, the maximum workday is seven hours, and they are not permitted to work in unhealthy or dangerous occupations or those requiring significant physical effort. Children must undergo a medical examination before beginning work. By law children must be paid at least the minimum wage or a minimum of two-thirds of the adult salary, whichever is higher.
Although the law prohibits forced and bonded labor by children, it was considered to be a common problem, especially in rural areas. Many children in rural areas were forced to work, particularly in commercial agriculture, as domestics, and in prostitution. The major factors contributing to the worst forms of child labor were chronic family poverty, lack of employment for adults, breakdown of family support mechanisms, the changing economic environment, lack of educational opportunities, gender inequality, and the impact of HIV/AIDS. Children, including those under the age of 15, commonly worked on family farms in seasonal harvests or on commercial plantations, where they picked cotton or tea leaves and were paid on a piecework basis.
The Ministry of Labor regulates child labor in both the informal and formal sectors. Labor inspectors may obtain court orders and use police to enforce compliance with child labor provisions. Violations of child labor provisions are punishable with fines ranging from one to 40 monthly salaries at minimum wage. Enforcement mechanisms generally were adequate in the formal sector but remained poor in the informal sector. The Labor Inspectorate and police forces lacked adequate staff, funds, and training to investigate child labor cases, especially in areas outside the capital where a majority of the abuses occurred. Although the government provided training for police on child prostitution and abuse, there was no specialized child labor training for the Labor Inspectorate. The government disseminated information and provided education about the dangers of child labor to the general public.
e. Acceptable Conditions of Work
In May the government granted a 14 percent increase in the statutory minimum wage for industry and services (including employees in public administration), bringing it to approximately $69 (1,645 meticais) per month. The government granted a 10 percent increase in the minimum wage in the agricultural sector bringing the monthly total to $47 (1,126 meticais). Despite the increase, which was slightly above the 8.2 percent inflation rate reported during the year, neither minimum wage provided a decent standard of living for a worker and family. Although the industrial sector frequently paid above minimum wage, there was little industry outside of the Maputo area. In addition less than 10 percent of workers held salaried positions, and the majority of the labor force worked in subsistence farming. Many workers used a variety of strategies to survive, including finding a second job, maintaining their own gardens, or depending on the income of other family members.
The Ministry of Labor is responsible for enforcing the minimum wage rates in the private sector and the Ministry of Finance in the public sector. Violations of minimum wage rates usually were investigated only after workers registered a complaint. Workers generally received benefits, such as transportation and food, in addition to wages.
The standard legal workweek is 40 hours but can be extended to 48 hours. After 48 hours, overtime must be paid at 50 percent over the base hourly salary. Overtime is limited by law to two hours per day and 100 hours per year. Foreign workers are protected under the law.
Worker complaints continued during the year concerning: employers deducting social security contributions from wages but failing to pay them into accounts; lack of access to the social security system; not adhering to the law concerning firings; and intimidation of union members.
In the small formal sector, health and environmental laws were in place to protect workers; however, the Ministry of Labor did not effectively enforce these laws, and the government only occasionally closed firms for noncompliance. There continued to be significant violations of labor laws in many companies and services. Workers have the right to remove themselves from work situations that endanger their health or safety without jeopardy to their continued employment; in practice threats of dismissal and peer pressure restricted this right.
In October an inspector from the Ministry of Labor found some 90 workers at the Golden Fields Flower Company (owned by former Foreign Minister Leonard Simao and his wife), in slave-like conditions, working long hours without proper protective equipment, living in tents, and with no access to sanitary facilities or safe drinking water. The workers had been recruited in Tete and Manica provinces, promised good working conditions, and provided transportation to Maputo. However, when the workers complained to the owners and asked to be provided transportation back to their home provinces, they were denied. Following the visit of the labor inspector, the government immediately suspended the company's operations and ordered the return of workers to their home provinces.
As of mid-September, the Ministry of Labor reported 62 labor accident victims, 40 of whom were temporarily incapacitated and 22 of whom were permanently incapacitated. While the law imposes fines for recurring accidents, no fines were imposed during the year. The law also requires that companies insure workers, but the Ministry of Labor estimates indicated that only between 50 and 60 percent of companies actually provided coverage. |
|Australian Journal of Educational Technology
1999, 15(3), 257-272.
This paper examines the relationship between the independent learner and computer based learning resources, which continue to be integral to educational delivery, especially in the training sector. To place interactivity in context, the first part of the discussion focuses on the major dimensions of interactivity and the different ways they have been characterised in computer based learning environments. These dimensions demonstrate the many ways that interactivity can be interpreted and the critical role that design and development plays in creating effective interactive encounters. The second part of the paper reviews the way storytelling structures and narrative have been promoted as effective strategies for enhancing comprehension and engagement in computer based learning applications. The way in which the interactivity and narrative are linked becomes critical to achieving this outcome.
Extending the use of a narrative within interactive media to include elements of performance and theatre, the third part of the discussion proposes that by conceptualising the learner as actor, a form of learner-designer communication can be established. Integrating this approach with elements of conversational and communication theory provides a context in which the learner-computer interface is transcended by that of learner and designer. Enabling this form of communication with the independent learner is suggested as a means to enhance computer based learning environments.
In contrast interactivity is portrayed as the distinguishing factor of the new media, with the assumption that "interactivity in a computer product means that the user, not the designer, controls the sequence, the pace, and most importantly, what to look at and what to ignore" (Kristof & Satran, 1995:35). This implies to some extent that the value of the material is based on the user's motivation and objectives. However, Holmes (1995:1-2) identifies that the structure or design of the material is equally critical and that the interactivity embodied within a computer is dependent on that design:
Interactivity is the ability of a new media program, web site, kiosk or multimedia presentation, to allow its user to control the content in some manner ... Interactivity must involve, engage, and motivate the user to explore the product ... Bad interactivity happens. The user can be frustrated by muddled organization, too much information or poor instructions ... Interactive properties of new media should provide opportunities for the user: exploration, discovery, and collaboration. Well conceived interactivity knows its audience, understands their knowledge base, and uses terms and phrases that are commonly understood by the audience. Good interactivity also takes into account the situation of the audience.Examining this perception reveals a number of major factors - first the user must be able to control the application and the application must be designed to engage the user; second the user must have some freedom and third the application must be designed for a specific audience. Interactivity, at least from this perspective, is partly about what the user does and part about how the application is designed. However, this perspective is at odds with the notion of the individual learner - are designers able to create applications that will match each member of a target group?
In proposing a nascent theory of interactivity Jaspers (1991) identified different levels ranging from the linear (information delivery) to the communicative (student initiated). However, the success of such a theory was dependent on the relationship between designer and developer and the perceived role of the learner. Emphasising the benefits of a learner centred environment the problematic and contradictory nature of interactivity was noted:
In fact, the expression of interactive delivery is contradicto in terminis from the viewpoint of interaction and emancipation. Delivery implies a unidirectional relationship. In full interaction there can be no one sided relation; the student is not just at the end of a chain but also at its beginning. (Jaspers (1991:21)To what extent the learner can play a proactive role while working with computer based learning resources becomes critical to the success of such applications and is considered in the final part of this paper. A related issue concerns developing a common understanding of interactivity. Schwier & Misanchuk (1993) refer to the levels defined in the literature as arbitrary and non-descriptive. Aldrich, Rogers & Scaife (1998) refer to interactivity as ubiquitous and Plowman (1996b) considers the excesses of physical interactions as gratuitous. The confused state of interactivity was also acknowledged in the following:
Though lauded by many for its ability to handle user inputs, there is little consensus with regard to the design of human computer interactions. Indeed, disagreement even exists about the meaning of the terms interactive as applied to emerging technologies. Researchers have described fundamentally different perspectives on the roles of interactions, ranging from facilitating lesson navigation to supporting encoding of specific lesson content ... Research on interaction methods may be among the most critical ... The domain of possibilities has broadened substantially, yet little research has been advanced which might guide their design. (Hannafin, Hannafin, Hooper, Rieber & Kini, 1996:385)A critical approach to assessing interactivity was presented by Rose (1999) who, given the various taxonomies of interactivity coupled with an apparent lack of denotative value, set about a deconstruction of established understandings. Critical of the "good" versus "bad" mentality that pervades the field (active not passive, learner control not program control, constructivism not instructivism, hypermedia not linear delivery), Rose (1999:45) presented an interesting scenario that contradicts many perceptions of interactive environments:
Texts addressing the subject of interactivity ... privilege representations of highly motivated learners exploring the wonderful worlds of interactive instructional programs and making exciting discoveries. However, it is by virtue of a deconstructive reading that we can begin to see the shadowy figure of the disavowed other lurking behind these wide eyed adventurers: the shadow of a child sitting mesmerised and immobile before the computer, only her index finger on the mouse moving occasionally as a stream of images passes in a more or less predetermined sequence before her eyes.Focusing explicitly on control Rose (1999) observed that although the words "learner control" have been interpreted as control by the learner, while grammatical comparisons (for example crowd control, weight control) tend to imply the opposite - control of the learner! The conclusion provides reinforcement for maintaining and extending research into making interactive learning work better:
That the field of educational computing is in need ... of internal critique is surely suggested by the fact that the very quality which is said to make computers unique and to justify their instructional use continues to defy definition. (Rose, 1999:48).There remains a paradox of interactivity. On the one hand it is portrayed as an integral and critical component of computer based applications and on the other as ill defined, deceptive and difficult if not impossible to implement. Is interactivity the hero or villain of achieving success through computer based learning? To address this question the following dimensions of interactivity demonstrate its multi-faceted nature and the issues that are essential to consistently implementing effective interactions within computer based structures.
However, the value of the control provided to the learner is dependent of the consequences of the learner actions and the extent to which the application responds or adapts to the individual learner's actions.
Interactive lessons are those in which the learner actively or overtly responds to information presented by the technology, which in turn adapts to the learner, a process more commonly referred to as feedback. The point is that interactive lessons require at least the appearance of two way communication. (Jonassen 1985:7)In this instance the adaptive capacity of the program is viewed in terms of a combination of the learner response and computer feedback, which in turn provides a form of communication. This process of question-response-feedback has been perceived by many writers as the essence of interactivity for computer based learning. For example, Steinberg (1991:100) observes that "question-response-feedback sequences help learners attain higher cognitive skills as well as factual information" and links the two elements of control and adaptation to the mechanics of interacting with the system (navigation) and the acquisition of knowledge and skills (learning). These descriptions identify two major components of interactivity - those which are program initiated and tell the learner what to do and those which are learner initiated, such as requests for help, information or explanation. This balance is also perceived as a means to achieve the goal of emulating human-human interaction in the computer medium, and thereby enhancing participation in the learning process.
Interactivity, or instructional features that promote active learning, provides critical support for increases in learning and retention in all educational activities ... Interaction implies active learner participation in the learning process ... an essential condition for effective learning ... failure to build interactivity into your program will reduce learning and retention.This comment reflects a shift in emphasis from the overt nature of interactivity to the extent to which internal learning is facilitated. In providing a set of guidelines for interactivity, Fenrich (1997) suggests the following set of options:
To provide a context for meaningful learning, Hannafin (1989) identified a set of five interactive functions (navigation, query, verification, elaboration, procedural control) and suggested a set of engaging activities to support these functions - fault free questions, queries, real time responding, notetaking, predicting/hypothesising, hypertext and cooperative dialogue. Although published over a decade ago the conclusions provide a useful guide for understanding interactivity:
It is no longer adequate to simply describe interactions in terms of either input technology employed or the physical characteristics of the responses made. ... We need a richer understanding of the psychological requirements associated with instructional tasks and responses, and a sense for how to extend design science beyond the methods that have evolved through the years. (Hannafin, 1989:178)Interactivity and meaningful learning has also been considered in terms of the schema model of human memory, where information is perceived as being stored in a web of interconnected nodes. "The strength of knowledge relies not simply on the number of nodes that exist, but more on the quality and number of interconnections between the nodes" Parrish (1996:2). Interactivity that encourages deeper cognitive processing (Craik & Lockhart, 1972) will potentially lead to these interconnections being built.
When considered in relation to computer enhanced learning, distance is an issue. While learner and computer are in close physical proximity, the extent to which there is distance (or lack of communication) between teacher, content and learner will diminish the effectiveness of the interactions. The extent to which communication or conversation between learner and designer can be integrated into a computer based medium has been analysed by Laurillard (1993), in which a conversational framework is used to identify the teacher-learner relationships. Enabling these concepts in the computer based medium is the challenge addressed within this paper.
Storytelling and narrative lie at the heart of all successful communication. Crude, explicit, button pushing interaction breaks the spell of engagement and makes it hard to present complex information that unfolds in careful sequence.In this scenario the difference between the overt interactions are explicitly compared with the potential of narrative. The challenge confronting educational technology developers is the production of computer based environments that engage the learner in effective instructional communication without the interactivity interfering with the overall process. Identifying the concepts of storytelling and narrative as critical determinants of communication have been shown to provide a context to enable the potential of interactive learning environments.
Since interactive interfaces ought to foster ... coordination between improvisation and planning we need to discover better theories of what is involved in the dynamic control of inquiry, line of thought, and action more generally. We need to discover more open ended models of coherence and narrative structure. Kirsh (1997:81)Interactivity is portrayed as the crucial element of the new technology and yet recent research has demonstrated that there is still much to understand about the ways in which the interactive process facilities access to technology, especially in the context of computer based learning applications. Plowman (1996a:102-103) states,
disruption of the narrative is strongest at the foci of interactivity ... (which) should be considered in terms of how they can be integrated into the overall narrative and how they can be used as a way of stimulating interest in the unfolding narrative ... by considering the interrelationship of narrative, linearity and interactivity and their design implications we can help learners to make sense of interactive multimedia.The extent to which the interactions implemented are disruptive or constructive is an area that has received little attention to date but one that may provide some new guidelines for designers. As a concept, narrative can be viewed as a linear storyline or in terms of how the story is told, the way it is received, what meanings it can have and the specific social, cultural, gendered and technological context in which it is told (Plowman, 1996a; Humphreys, 1997b). But in what ways might narrative assist meaning, reduce the impact of interactive interference and provide the necessary framework to promote learning amongst diverse groups of learners? Plowman (1996a:92) suggests that:
Narrative coherence is identified here with a lack of redundancy and a fixed sequence. Interactive multimedia (IMM) programmes challenge these traditional definitions of narrative because it can be suspended or altered at discrete decision points, the foci of interactivity, and a rearrangement of discrete elements gives rise to new text and new meanings. While the concepts of wholeness, unity and coherence of meanings are not fashionable in a post-modern world, in educational multimedia... the notion of multiple interpretations has different implications, particularly for comprehension and cognition.While "narrative isn't just a shaping device: it helps us think, remember, communicate, and make sense of ourselves and the world" (Plowman, 1996a:3), when considered in terms of interactive environments, its perceived advantages are through enhancing comprehension and understanding:
Narrative structure is fundamental to comprehension to the extent that when it is clearly absent from certain forms of multimedia (it) can seriously undermine comprehension of the material. (Laurillard, 1998:231).
Framework which accommodates the audience or user into the process of engagement with interactive media in ways that narrative theory finds difficult... Interactivity produces for the user of media a different relationship to story. This shift in relationship may be able to be framed as a shift from narrative, as an experience of recounting a story, to play, as an experience of enacting a story. Humpherys (1997:9,11)A focus on either narrative or play as structures to enhance engagement implies that the interactivity provided must be integrated to enable reinforcement of the specific learning objectives and to maintain the user's participation in the story. It is our challenge therefore to develop applications that minimise the potential for interactive interference. By considering the concept of interactivity, narrative and play in association with the links between the designer and learner, the following section presents a model to potentially enhance the communication between learner and designer within computer based environments.
The association of education with theatre is by no means novel:
As we develop as persons, we develop a sophistication and sensitivity to what is the "proper" role behaviour for various groups we must meet ... We are in a sense able to predict the consequences of various behaviour alternatives on others without actually performing them, and can select the best role and performance ... This process is seen as "dramatic rehearsal" ... and is as complex as the relationship of actor and audience on stage. Hodgkinson (1967:3)In the context of education and social change, these concepts of performance, roles, identity and cues provide a focus for the learning environment. "As on stage, we are constantly sending out signals to those around us telling them how we wish to behave ... education at all levels is constantly faced with the problem of correcting misunderstood or unintentional cues" (Hodgkinson, 1967:22). To what extent therefore do computer based learning applications provide learners with a set of confusing cues, and in what way are learners allowed to play a role when working with educational technology? If, as Hodgkinson (1967) argues, being unaware of role playing can damage educational goals, it may be that the concept of roles, cues, theatre and performance hold clues as to making computer based learning work better.
Laurel (1990), who proposed the computer as theatre, also identifies the importance of human agency in human-computer interaction. Using this analogy, it is possible to extend the narrative and play concept to that of theatre and the role of the learner from observer or participant to actor:
By explicitly casting the language learner as actor (or other), a more playful and reflexive context for taking performative risks becomes possible. At the same time, the learner is pressed to assume responsibility for communicative acts that involve skill building at multiple levels of performance (phonological, kinetic, pragmatic), and that include but go beyond propositional knowledge. (Quinn, 1997:1).If the learner is to become and actor, then the structure of the application can be perceived as a performance. Trognazzini (1999) identifies magicians and illusionists as performers and that, in the same way that their act must be believable, so must interactive software have the same elements of acceptance:
I propose that there is a "threshold of believability," a point at which careful design and meticulous attention to detail have been sufficient to arouse in the spectator or user a belief that the illusion is real. The exact point will vary by person and even by mood, so we must exceed it sufficiently to ensure believability.If the learner believes the illusion created by the educational environment then the interactions will become integral to that illusion rather than as external controls to an environment being observed. Playing a role, being an actor in the learning process is not only about making choices but becoming part of the narrative, story or performance.
To date, many computer based learning resources have placed the learner as observer (passive or active) in much the same way as members of an audience at a performance. However, this separation of actor and audience provides an essentially uni-directional communication - from stage to audience. If the learner were to take a position on stage, what might that mean for computer based learning and how might it be implemented?
Designers might extend the concepts proposed by Hannafin & Peck (1988) to create scenarios in which the learner has an opportunity to participate. This would involve the learner being asked how familiar they are with the application and the extent to which they would like orientation to the application. The tour should allow the user to ask questions for clarification, in much the same way that an actor and director peruse the script and work through a series of rehearsals. Once comfortable with the location and prepared for "opening night", learners need to be made familiar with the controls (the stage, props and other actors) - not in their use but in their purpose - and the relevance of their appearance on the display (location on the stage). By using a narrative or story to define the performance in which the learner is participating, a logical and meaningful series of interactions can be employed.
An argument against such design structures is that they are too expensive or too difficult to implement, but this is not necessarily the case. If computer based learning resources continue to be used then it is important that they be effective. To date this has not been achieved consistently, so if more effective resource can be structured then the initial costs will be outweighed by the benefits. Similarly, as technology is developing so rapidly the creation of what now appear to be complex environments will become components of the development software.
The integration of the learner into the overall process may provide an environment in which the communication is focused on learner and teacher (designer) rather than learner and computer (content). Laurillard (1993) proposed the discursive, adaptive, interactive and reflective elements of education media in relation to effective teacher-learner communications. Continuing to develop the way in which people work with computers in learning contexts will provide the means by which these elements will continue to be successfully integrated into computer based environments. Conceptualising the learner as actor may provide a means to achieve this.
One of the first mottos I heard when beginning my career in this field was "you're only limited by your imagination". My vision is one of a learner integrated, engaged and achieving in a computer based environment - imagining a learner on stage, playing a critical role in the narrative is how that vision might be realised. The success of computer based learning will be through interactivity as a manifestation of communication between designer and learner. If the designer can develop their ideas into a performance into which the learner is actor and interactivity the stage, then the illusion, magic and engagement so eagerly sought after might well be achieved.
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|Author: Rod Sims, School of Multimedia and Information Technology, Southern Cross University. Email: email@example.com
Please cite as: Sims, R. (1999). Interactivity on stage: Strategies for learner-designer communication. Australian Journal of Educational Technology, 15(3), 257-272. http://www.ascilite.org.au/ajet/ajet15/sims.html |
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum IV. Materials R&D—A Vision of the Scientific Frontier The Science of Modern Technology Paul Peercy SEMI/SEMATECH Electronic, Optical, and Magnetic Materials and Phenomena We have seen numerous important unexpected discoveries in all areas of condensed-matter and materials physics in the decade since Physics Through the Nineties was published. Although these scientific discoveries are extremely impressive, perhaps equally impressive are the technological advances based on our ever-increasing understanding of the basic physics of materials along with our increasing ability to tailor the composition and structure of materials in a cost-effective manner. Today's technological revolution would not be possible without the continuing increase in our scientific understanding of materials and phenomena, along with the processing and synthesis required for high-volume, low-cost manufacturing. This article examines selected examples of the scientific and technological impact of electronic, optical, and magnetic materials and phenomena. Technology based on electronic, optical, and magnetic materials is driving the information age through revolutions in computing and communications. With the miniaturization made possible by the invention of the transistor and the integrated circuit (IC), enormous computing and communication capabilities are becoming readily available worldwide. These technological capabilities enabled the Information Age and are fundamentally changing how we live, interact, and transact business. These technologies provide an excellent demonstration of the strong interdependence and interplay of science and technology. They have greatly expanded the tools and capabilities available to scientists and engineers in all areas of research and development, ranging from basic physics and materials research to other areas of physics and to such diverse fields as medicine and biotechnology. Incorporation of major scientific and technological advances into new products can take decades and often follows unpredictable paths. Selected technologies supported by the foundations of electronic, photonic, and magnetic phenomena and materials are illustrated in Figure 1 . These technologies have enabled breakthrough technologies in virtually every sector of the national economy. The two-way interplay between foundations and technology is a major driving force in this field. The most recent fundamental advances and technological discoveries have yet to realize their potential. Figure 1. Examples of how major scientific and technological advances have an impact on new products. The Science of Information Age Technology The predominant semiconductor technology today is the silicon-based integrated circuit. The silicon integrated circuit is the engine that drives the information revolution. For the past 30 years, the technology has been dominated by Moore's law—the statement that the density of transistors on a silicon integrated circuit doubles about every 18 months. The relentless reduction in transistor size and increase in circuit density have provided the increased functionality per unit cost that underlies the information revolution. Today's computing and communications capability would not be possible without the phenomenal 25 to 30 percent per year exponential growth in capability per unit cost since the introduction of the integrated circuit in about NOTE: This article was prepared from written material provided to the Solid State Sciences Committee by the speaker.
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum 1960. That sustained rate of progress has resulted in low-cost volume manufacturing of high-density memories with 64 million bits of memory on a chip and complex, high-performance logic chips with ~10 million transistors on a chip. This trend is projected to continue for the next several years. If the silicon integrated circuit is the engine that powers the computing and communications revolution, optical fibers are the highways for the Information Age. Although fiber optics is a relatively recent entrant in the high technology arena, the impact of this technology is enormous and growing. It is now the preferred technology for transmission of information over long distances. There are already approximately 30 million km of fiber installed in the United States and an estimated 100 million km installed worldwide. Due in part to the faster than exponential growth of connections to the Internet, the installation of optical fiber worldwide is occurring at an accelerated rate of over 20 million km per year—more than 2,000 km/h, or around Mach 2. In addition, the rate of information transmission down a single fiber is increasing exponentially at a rate of a factor of 100 every decade. Transmission in excess of 1 terabit per second has been demonstrated in the research laboratory, and the time lag between laboratory demonstration and commercial system deployment is about 5 years. Compound semiconductor diode lasers provide the laser photons that are the vehicles that transport information along the optical information highways. Semiconductor diode lasers are also at the heart of optical storage and compact disk technology. In addition to their use in very-high-performance microelectronics applications, compound semiconductors have proven to be an extremely fertile field for advancing our understanding of fundamental physical phenomena. Exploiting decades of basic research, we are now beginning to be able to understand and control all aspects of compound semiconductor structures, from mechanical through electronic to optical, and to grow devices and structures with atomic layer control, in a few specific materials systems. This capability allows the manufacture of high-performance, high-reliability, compound semiconductor diode lasers that can be modulated at gigahertz frequencies to send information over the fiber-optical networks. High-speed semiconductor-based detectors receive and decode this information. These same materials provide the billions of light-emitting diodes sold annually for displays, free-space or short-range high-speed communication, and other applications. In addition, very-high-speed, low-power compound semiconductor electronics play a major role in wireless communication, especially for portable units and satellite systems. Another key enabler of the information revolution is low-cost, low-power, high-density information storage that keeps pace with the exponential growth of computing and communication capability. Both magnetic and optical storage are in wide use. Very recently, the highest-performance magnetic storage/readout devices have begun to rely on giant magnetoresistance (GMR), a phenomenon that was discovered by building on more than a century of research in magnetic materials. Although Lord Kelvin discovered magnetoresistance in 1856, it was not until the early 1990s that commercial products using this technology were introduced. In the last decade, the condensed-matter and materials understanding converged with advances in our ability to deposit materials with atomic-level control to produce the GMR heads that were introduced in workstations in late 1997. It is hoped that, with additional research and development, spin valve and colossal magnetoresistance technology may be understood and applied to workstations of the future. This increased understanding, provided in part by our increased computational ability arising from the increasing power of silicon ICs, coupled with atomic-level control of materials, led to exponential growth in the storage density of magnetic materials analogous to Moore' s law for transistor density in silicon ICs. Future Directions and Research Priorities Numerous outstanding scientific and technological research needs have been identified in electronic, photonic, and magnetic materials and phenomena. If those needs are met, it is anticipated that these technology areas will continue to follow their historical exponential growth in capability per unit cost for the next few years. Silicon integrated circuits are expected to follow Moore's law at least until the limits of optical lithography are reached, transmission bandwidth of optical fibers is expected to grow exponentially with advances in optical technology and the development of soliton propagation, and storage density in magnetic media is expected to grow exponentially with the maturation of GMR and development of colossal magnetoresistance in the not too distant future. Although these changes will have a major impact on computing and communications over the next few years, it is clear that extensive research will be required to produce new concepts and that new approaches must be developed
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum to reduce research concepts to practice if these industries are to maintain their historical growth rate over the long term. Continued research is needed to advance the fundamental understanding of materials and phenomena in all areas. More than a century of research in magnetic materials and phenomena has given us an understanding of many aspects of magnetism, but we still lack a comprehensive first-principles understanding of magnetism. By comparison, the technology underlying optical communication is very young. The past few years have seen enormous scientific and technological advances in optical structures, devices, and systems. New concepts such as photonic lattices, which are expected to have significant technological impact, are emerging. We have every reason to believe that this field will continue to advance rapidly with commensurate impact on communications and computing. As device and feature sizes continue to shrink in integrated circuits, scaling will encounter fundamental physical limits. The feature sizes at which these limits will be encountered and their implications are not understood. Extensive research is needed to develop interconnect technologies that go beyond normal metal and dielectrics in the relatively near term. Longer term, technologies are needed to replace today 's Si field-effect transistors. One approach that bears investigation is quantum state switching and logic as devices and structures move further into the quantum mechanical regime. A major future direction is nanostructures and artificially structured materials, which was a general theme in all three areas. In all cases, artificially structured materials with properties not available in nature revealed unexpected new scientific phenomena and led to important technological applications. As sizes continue to decrease, new synthesis and processing technologies will be required. A particularly promising area is self-assembled materials. We need to expand research in self-assembled materials to address such questions as how to controllably create the desired one-, two-, and three-dimensional structures. As our scientific understanding increases and synthesis and processing of organic materials systems mature, these materials are expected to increase in importance for optoelectronic, and perhaps electronic, applications. Many of the recent technological advances are the result of strong interdisciplinary efforts as research results from complementary fields are harvested at the interface between the fields. This is expected to be the case for organic materials; increased interdisciplinary efforts, for example between CMMP, chemistry, and biology, offer the promise of equally impressive advances in biotechnology. Conclusion In conclusion, we identify a few major scientific and technological questions that are still outstanding and call attention to research and development priorities. Selected Major Unresolved Scientific and Technology Questions What technology will replace normal metals and dielectrics for interconnect in silicon ICs as speed continues to increase? What is beyond today's field-effect transistorbased Si technology? Can we create an all-optical communications/computing network Can we understand magnetism on the mesoscales and nanoscales needed to continue to advance technology? Can we fabricate devices with 100 percent spinpolarized current injection? Priorities Advance synthesis and processing techniques, including nanostructures and self-assembled one-, two-, and three-dimensional structures; Pursue quantum state logic; Exploit physics and materials science for low-cost manufacturing; Pursue the physics and chemistry of organic and other complex materials for optical, electrical, and magnetic applications; Develop techniques to magnetically detect individual electron and nuclear spins with atomic-scale resolution; and Increase partnerships and cross-education/communications among industry, university, and government laboratories.
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum Novel Quantum Phenomena in Condensed-Matter Systems Steven M. Girvin Indiana University The various quantum Hall effects (QHEs) are arguably some of the most remarkable many-body phenomena discovered in the second half of the 20th century, comparable in intellectual importance to superconductivity and superfluidity. They are an extremely rich set of phenomena with deep and truly fundamental theoretical implications. The fractional effect, for which the 1998 Nobel Prize in Physics was awarded, has yielded fractional charge, spin, and statistics, as well as unprecedented order parameters. There are beautiful connections with a variety of different topological and conformal field theories studied as formal models in particle theory, each here made manifest by the twist of an experimental knob. Where else but in condensedmatter physics can an experimentalist change the number of flavors of relativistic chiral Fermions or set by hand the Chem-Simons coupling that controls the mixing angle for charge and flux in 2+1D electrodynamics? Because of recent technological advances in molecular beam epitaxy and the fabrication of artificial structures, the field continues to advance with new discoveries even well into the second decade of its existence. Experiments in the field were limited for many years to simple transport measurements that indirectly determine charge gaps. However recent advances have led to many successful new optical, acoustic, microwave, specific heat, and nuclear magnetic resonance (NMR) probes, which continue to advance our knowledge as well as raise intriguing new puzzles. The QHE takes place in a two-dimensional electron gas subjected to a high magnetic field. In essence, it is a result of commensuration between the number of electrons, N, and the number of flux quanta, NΦ, in the applied magnetic field. The electrons undergo a series of condensations into new states with highly nontrivial properties whenever the filling factor ν = N/NΦ takes on simple rational values. The original experimental manifestation of the effect was the observation of an energy gap yielding dissipationless transport (at zero temperature) much like in a superconductor. The Hall conductivity in this dissipationless state is universal, given by σxy = ve2/h independent of microscopic details. As a result of this, it is possible to make a high-precision determination of the fine structure constant and to realize a highly reproducible quantum mechanical unit of electrical resistance, now used by standards laboratories around the world to maintain the ohm. The integer quantum Hall effect owes its origin to an excitation gap associated with the discrete kinetic energy levels (Landau levels) in a magnetic field. The fractional quantum Hall effect has its origins in very different physics of strong Coulomb correlations, which produce a Mott-insulator-like excitation gap. In some ways, however, this gap is more like that in a super-conductor, because it is not tied to a periodic lattice potential. This permits uniform charge flow of the incompressible electron liquid and hence a quantized Hall conductivity. The microscopic correlations leading to the excitation gap are captured in a revolutionary wave function developed by R.B. Laughlin that describes an incompressible quantum liquid. The charged quasi particle excitations in this system are “anyons” carrying fractional statistics intermediate between bosons and Fermions and carrying fractional charge. This sharp fractional charge, which despite its bizarre nature has always been on solid theoretical ground, has recently been directly observed two different ways. The first is an equilibrium thermodynamic measurement using an ultrasensitive electrometer built from quantum dots. The second is a dynamical measurement using exquisitely sensitive detection of the shot noise for quasi particles tunneling across a quantum Hall device. Quantum mechanics allows for the possibility of fractional average charge in both a trivial way and a highly nontrivial way. As an example of the former, consider a system of three protons forming an equilateral triangle and one electron tunneling among the 1S atomic bound states on the different protons. The electronic ground state is a symmetric linear superposition of quantum amplitudes to be in each of the three different 1S orbitals. In this trivial case, the mean electron number for a given orbital is 1/3. This, however, is a result of statistical fluctuations because a measurement will yield electron number 0 two-thirds of the time and electron number 1 one-third of the time. These fluctuations occur on a very slow time scale and are associated with the fact that the electronic spectrum NOTE: This article was prepared from written material provided to the Solid State Sciences Committee by the speaker.
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum consists of three very nearly degenerate states corresponding to the different orthogonal combinations of the three atomic orbitals. The ν = 1/3 QHE has charge 1/3 quasi particles but is profoundly different from the trivial scenario just described. An electron added to a ν = 1/3 system breaks up into three charge 1/3 quasi particles. If the locations of the quasi particles are pinned by (say) an impurity potential, the excitation gap still remains robust and the resulting ground state is nondegenerate. This means that a quasi particle is not a place (like the proton above) where an extra electron spends one-third of its time. The lack of degeneracy implies that the location of the quasi particle completely specifies the state of the system, that is, implies that these are fundamental elementary particles with charge 1/3. Because there is a finite gap, this charge is a sharp quantum observable that does not fluctuate (for frequencies below the gap scale). The message here is that the charge of the quasi particles is sharp to the observers as long as the gap energy scale is considered large. If the gap were 10 GeV instead of 10 K, we (living at room temperature) would have no trouble accepting the concept of fractional charge. Magnetic Order of Spins and Pseudospins At certain filling factors (ν = 1, in particular) quantum Hall systems exhibit spontaneous magnetic order. For reasons peculiar to the band structure of the GaAs host semiconductor, the external magnetic field couples exceptionally strongly to the orbital motion (giving a large Landau level splitting) and exceptionally weakly to the spin degrees of freedom (giving a very small Zeeman gap). The resulting low-energy spin degrees of freedom of this ferromagnet have some rather novel properties that have recently begun to be probed by NMR, specific heat, and other measurements. Because the lowest spin state of the lowest Landau is completely filled at ν = 1, the only way to add charge is with reversed spin. However, because the exchange energy is large and prefers locally parallel spins (and because the Zeeman energy is small), it is cheaper to partially turn over several spins forming a smooth topological spin “texture.” Because this is an itinerant magnet with a quantized Hall conductivity, it turns out that this texture (called a skyrmion by analogy with the corresponding object in the Skyrme model of nuclear physics) accommodates precisely one extra unit of charge. NMR Knight shift measurements have confirmed the prediction that each charge added (or removed) from the ν = 1 state flips over several (~ 4 to 30 depending on the pressure) spins. In the presence of skyrmions, the ferromagnetic order is no longer collinear, leading to the possibility of additional low-energy spin wave modes, which remain gapless even in the presence of the Zeeman field (somewhat analogous to an antiferromagnet). These low-frequency spin fluctuations have been indirectly observed through a dramatic enhancement of the nuclear spin relaxation rate 1/T1. In fact, under some conditions T1 becomes so short that the nuclei come into thermal equilibrium with the lattice via interactions with the inversion layer electrons. This has recently been observed experimentally through an enormous enhancement of the specific heat by more than five orders of magnitude. Spin is not the only internal degree of freedom that can spontaneously order. There has been considerable recent progress experimentally in overcoming technical difficulties in the MBE fabrication of high-quality multiple-well systems. It is now possible for example to make a pair of identical electron gases in quantum wells separated by a distance (~ 100 Å) comparable to the electron spacing within a single quantum well. Under these conditions, strong interlayer correlations can be expected. One of the peculiarities of quantum mechanics is that, even in the absence of tunneling between the layers, it is possible for the electrons to be in a coherent state in which their layer index is uncertain. To understand the implications of this, we can define a pseudospin that is up if the electron is in the first layer and down if it is in the second. Spontaneous interlayer coherence corresponds to spontaneous pseudospin magnetization lying in the XY plane (corresponding to a coherent mixture of pseudospin up and down). If the total filling factor for the two layers is ν = 1, then the Coulomb exchange energy will strongly favor this magnetic order just as it does for real spins as discussed above. This long-range transverse order has been observed experimentally through the strong response of the system to a weak magnetic field applied in the plane of the electron gases in the presence of weak tunneling between the layers. Another interesting aspect of two-layer systems is that, despite their extreme proximity, it is possible to make separate electrical contact to each layer and perform drag experiments in which current in one layer induces a voltage in the other due to Coulomb or phonon-mediated interactions. Stacking together many quantum wells gives an artificial three-dimensional structure analogous to certain
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum organic Bechgaard salts in which the QHE has been observed. There is recent growing interest in the bulk and edge (“surface”) states of such three-dimensional systems and with the nature of possible Anderson localization transitions. These phenomena and numerous others, which cannot be mentioned because of space limitations, have provided a wonderful testing ground for our understanding of strongly correlated quantum ground states that do not fit into the old framework of Landau's Fermi liquid picture. As such, they are providing valuable hints on how to think about other strongly correlated systems such as heavy Fermion materials and high-temperature superconductors. Nonequilibrium Physics James S. Langer University of California, Santa Barbara Nonequilibrium physics is concerned with systems that are not in mechanical or thermal equilibrium with their surroundings. Examples include flowing fluids under pressure gradients, solids deforming or fracturing under external stresses, and quantum systems driven by magnetic fields. These systems often lead to very familiar patterns such as snowflakes, dendritic microstructures in alloys, or chaotic motions in turbulent fluids. Many of these are familiar phenomena governed by well-understood equations of motion (e.g., the Navier Stokes equation), but in some of the most interesting cases, the implications of these equations are not understood. The Brinkman report (Physics Through the 1990s, National Academy Press, Washington, D.C., 1986) recognized the significance of the emerging field of nonequilibrium physics but missed some of the most important topics of current research such as friction, fracture, and granular materials. Notable progress has been made in the last decade regarding patterns in convecting and vibrating fluids, reaction-diffusion systems, aggregation, and membrane morphology. The patterns observed in nonequilibrium systems are especially sensitive to small perturbations. Weather phenomena are a prime example. Long-range weather forecasting requires precise characterization of current and past weather conditions. As such characterization becomes more detailed, it is possible to predict future patterns with increasing accuracy. One task of nonequilibrium physics is quantifying the relationship between precision and predictability. In spite of decades of study, the origin of ductility in materials remains a key unsolved problem of nonequilibrium physics. Traditional explanations based on dislocations do not explain observations such as ductility in glassy materials. We lack a good theory of ductile yielding in situations where stresses and strains vary rapidly in space and time. One of the most important recent observations is that fast brittle cracks undergo materials-specific instabilities leading to roughness on the fracture surface. Stick-slip friction is also observed on large scales, for example, in earthquake dynamics. The nature of these processes, including the issue of lubricated friction, is a key problem of nonequilibrium physics. Granular materials are an example of a familiar class of materials of considerable industrial importance that have escaped scrutiny by physicists until recently. These materials are highly inelastic in their interactions. When granular materials cohere slightly, they can behave like viscous fluids as in saturated soil. If the coherence is strong, then we have sandstone or concrete, which behaves more like ordinary solids. If complex dynamics are added, we have foams or dense colloidal suspensions. The nature of lubrication is also relevant to these problems. These and many other open problems show that the frontiers of physics include many very familiar phenomena. In many cases, these problems are of great importance to materials properties and industrial processes.
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum Soft Condensed Matter V. Adrian Parsegian National Institutes of Health Adrian Parsegian opened his remarks on soft condensed matter research with the question, “When was American poetry born?” He quoted William Saroyan's response, “When people not trained in poetry began writing it.” No one can imagine Walt Whitman's poetry being written in England or anywhere except America. It was “an instantaneous flop” because people did not think it was even “poetry.” The clear implication for the field of condensed-matter and materials physics is that soft materials physics will flourish only after much initial skepticism and even resistance are overcome. Advancing our understanding of soft materials requires an unprecedented combination of traditionally distinct, and noninteracting, fields. One of the greatest challenges is the simple recognition and appreciation of the disparate skills required for attacking problems with relentlessly increasing levels of complexity. For example, as the genome project unfolds, uncovering seemingly endless genetic information, synthetic chemists and physicists face the daunting task of producing and understanding the complex interactions that govern biological function occurring at length scales ranging from atomic to supramolecular dimensions. Such problems may not succumb to the conventional reductire methods familiar to most physicists. Modern instrumentation such as third-and even fourth-generation synchrotrons, high-flux neutron sources, and high-resolution nuclear magnetic resonance spectrometers provide powerful means for exploring these issues. However, are these potent tools the key to uncovering the secrets of biology? “What constitutes understanding in this business?” asks Parsegian, adding, “An explanation that satisfies the physicist may be thoroughly irrelevant to the gene therapist.” Established approaches toward education and funding and even attitudes about industrial interactions must change if the physics community is to have a demonstrable and meaningful impact on this burgeoning field. Many of the macroscopic properties of soft materials are foreign to the condensed-matter physicist. Softness is substituted for hardness as a desirable property; malleability, extensibility, and compliance replace stiffness and shape retention; and fragility is often more valuable than durability. These themes are stimulating a new type of physics that relies on the same bedrock principles enumerated in elementary physics education but must be augmented by the targeted interdisciplinary studies of medicine, food, polymers, and many others. The field of polymers provides a bridge between conventional physics and the biologically oriented sciences and engineering. Both natural and synthetic macromolecules offer numerous research and development opportunities. Polysaccharides and milk proteins are identified by Parsegian as examples of ubiquitous and naturally occurring macromolecules that may be formulated into novel items of commerce for use in foods and environmentally benign—even biodegradable—plastics. Advances in synthetic chemistry during the past decade have greatly expanded our ability to tailor molecular architecture, even in the simplest of polymers known as polyolefins, prepared from just carbon and hydrogen. This class of synthetic polymers makes up roughly 60 percent of the entire synthetic polymer market. Yet the commercial consequences of varying the number and length of side branches, grafts, and block sequencing on the melt flow properties, crystallization kinetics, and ultimate mechanical properties are just now being realized. Dendrimers, precisely and highly branched giant molecules that can assume a nearly perfect spherical topology, offer fresh strategies for manipulating polymer rheology and may provide ideal substrates for delivering drugs to the human body. Parsegian noted that polyelectrolytes are an especially important class of materials, since almost all forms of biologically relevant matter contain macro-molecules with some degree of ionic charging. Yet polyelectrolytes present some of the toughest challenges to condensed matter physics, convoluting electrostatic interactions, self-avoiding chain statistics, and traditional solution thermodynamics with self-assembly into higher-order structures that rely on tertiary and quaternary interactions. Application of physics to biological problems is not a new phenomenon. Physicists have made significant contributions to the field of protein folding for nearly 35 years. In fact, physicists have defined the “language” of the field and created exquisite tools for simulating and even “watching” individual molecules. For example, optical tweezers techniques permit quantification of the force versus extension relationship of indi-
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum vidual biomolecules such as DNA. These studies are not yet the clinical analysis of protein-folding and prionrelated illnesses such as Alzheimer 's and “Mad Cow” disease. “Why are physicists frequently off the biological radar screen? “ asks Parsegian. “In large part they are seen as insular and parochial,” is the response he has received to this question. He points out key differences in how physicists approach problemsolving through a simple example—understanding the force required to separate two interacting biomolecules, as might be encountered in cell adhesion. This delicate problem depends on the time scale of the experiment as much as the force applied and the displacement measured. Given enough time, the molecules will disengage without any applied force. Physicists must resist the temptation to connect or reconfigure anything biological into a conventional, solved, physics problem, in this case treating the biomolecular interactions with traditional static intermolecular potentials. Rectifying these shortcomings, that is, making physics visible and relevant to the biological sciences, will require education on the part of biologists as well as physicists. Biologists and medical researchers must understand the utility and importance of physics to their work. For example, the sophisticated instrumentation that is often the source of so many scientific revelations (e.g., three-dimensional NMR and x-ray imaging) often comes from fundamental advances in physics. Medical doctors and researchers should understand the origins of their equipment and appreciate the underlying principles of operation. Parsegian offered an assortment of recommendations for improving the current situation. Basic research, which fuels practical developments in industry and medicine, would benefit from the following innovations: (1) grant mechanisms that encourage interdisciplinary work; (2) special grants that circumvent the double jeopardy of being judged both as a biologist and a physicist; (3) fellowships for physicists, including theorists, to work in biological laboratories; and (4) maximized contact between university researchers and industrial scientists, especially those from the chemical, medical, and pharmaceutical industries. Changes in education were also prescribed. New physics courses must be developed that are targeted at biologists; and physicists should be trained in chemistry, biochemistry, and molecular biology. Introductory physics courses must begin to emphasize soft systems in addition to the traditional curriculum. There should be “bilingual” textbooks, aimed at both physicists and biologists. Summer schools with laboratories for scientists at all stages of their careers and interdisciplinary workshops can be established. In short, the field of physics should spread itself out from the confines of physics departments, while broadening its horizons to encompass the emerging exciting world of soft matter. Fractional Charges and Other Tales from Flatland Horst Störmer Bell Laboratories and Columbia University Flatland is two-dimensional space. In nature it is found at surfaces or at interfaces. The quantum Hall effect (QHE) and fractional quantum Hall effect (FQHE) are properties of electrons confined to the interface region of semiconductor quantum wells. The electrons can move along the two-dimensional surface of the interface but are confined in the third direction. A Nobel prize was awarded to Klaus von Klitzing in 1985 for the discovery of the QHE, and Störmer, D.C. Tsui, and Robert Laughlin shared the 1998 Nobel Prize in Physics for their discovery of the FQHE. This research area continues to be interesting, with many new ideas and discoveries. E.H. Hall discovered in 1878 a transverse voltage Vxy when a magnetic field B is imposed perpendicular (z-direction) to the direction of electrical current. The voltage is proportional to the current Ix in the layer. The ratio Vxy/Ix = Rxy defines the transverse resistance. The QHE is the observation of plateaus in the transverse resistance when measured as a function of magnetic field. These plateaus occur when Landau levels are completely filled with electrons. During these plateaus, the longitudinal resistance (RXX = VXX/Ix) appears to vanish: It actually declines to a very small value. The plateaus have a value of resistance that are multiples of the fundamental value h/e2 = 258120, where h is Planck's constant and e is the unit of charge. The plateaus have the same value in each sample. They have become the new international standard of resistance. Störmer, Tsui, and Gossard continued the measurements of the Hall effect to very high values of mag-
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Materials in a New Era: Proceedings of the 1999 Solid State Sciences Committee Forum netic field. They discovered additional plateaus in the transverse resistance. These plateaus occur at values that imply that the Landau levels are fractionally occupied: 1/3, 2/5, and so on. The experiments continue, with samples of increasing purity, at lower temperature, and with higher magnetic field. Low-temperature values of the mobility in GaAs/AlGaAs wells are now μ = 2 × 107 cm2/(Vs). The mean-free-path of the electrons is on the order of 100 × 10−6 m. In these recent measurements, the list of fractions is quite large, with unlikely numbers such as 5/23 appearing. Why are there fractions? The two-dimensional gas of electrons becomes highly correlated. It is difficult to imagine a many-body system that is as simple and as profound. At these high values of magnetic field, the electrons have circular orbits, due to the field, which have a small radius. The separation between electrons is much longer then the diameter of the orbits. The electrons are becoming localized into a kind of Wigner crystal. However, they have unusual motion, because they have no kinetics. Their correlation is entirely due to their mutual Coulomb interaction. The fractions indicate the existence of highly correlated states that occur at these fractional fillings. The topic is fascinating to theoretical physicists, who try to explain the origin of all of these states. The most important fraction is 1/3. It was discovered first and has a large plateau in the Hall resistance. Each quantum of flux can be considered to be a vortex in the plane. For the 1/3 state, there are three vortices for each electron. The three vortices get attached to the electron and form a quasi particle called a “composite.” For the 1/3 state, it is a “composite Boson,” and the collective state is due to a Bose-Einstein condensation of these quasi particles. Noise measurements in the resistivity confirm that the charge on the current carriers is actually e/3. The excitation energy to excite an electron out of the correlated state is ∆E = 10 K in temperature units. This energy is large, because the experiments are performed at a small fraction of a degree. At the fraction 1/2, the quasi particles are “composite Fermions,” which cannot form a condensate. There are no plateaus at this fraction, although it is speculated that the Fermions form a superfluid state akin to the Bardeen-Cooper-Schrieffer state in a superconductor. This pairing of Fermions is thought to explain the features of the 5/2 state in particular. For fractions written as the ratio of two integers q/p, even values of p are composite Fermions and odd values of p denote composite Bosons. New fractions continue to be discovered. The graph of the transverse resistance appears to be a “devil's staircase” of an infinite number of steps of irregular width. Physicists have always been fascinated by the highly correlated electron gas. There is no more highly correlated system of electrons than is found in the FQHE.
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Representative terms from entire chapter: |
BOOK I: INTRODUCTION
View of the Aztec Civilization
ANCIENT MEXICO- ITS CLIMATE AND ITS PRODUCTS- ITS PRIMITIVE RACES- AZTEC EMPIRE
THE country of the ancient Mexicans, or Aztecs as they were called, formed but a very small part of the extensive territories comprehended in the modern republic of Mexico. Its boundaries cannot be defined with certainty. They were much enlarged in the latter days of the empire, when they may be considered as reaching from about the eighteenth degree north to the twenty-first on the Atlantic; and from the fourteenth to the nineteenth, including a very narrow strip, on the Pacific. In its greatest breadth, it could not exceed five degrees and a half, dwindling, as it approached its south-eastern limits, to less than two. It covered, probably, less than sixteen thousand square leagues. Yet, such is the remarkable formation of this country, that though not more than twice as large as New England, it presented every variety of climate, and was capable of yielding nearly every fruit found between the equator and the Arctic circle.
All along the Atlantic the country is bordered by a broad tract, called the tierra caliente, or hot region, which has the usual high temperature of equinoctial lands. Parched and sandy plains are intermingled with others of exuberant fertility, almost impervious from thickets of aromatic shrubs and wild flowers, in the midst of which tower up trees of that magnificent growth which is found only within the tropics. In this wilderness of sweets lurks the fatal malaria, engendered, probably, by the decomposition of rank vegetable substances in a hot and humid soil. The season of the bilious fever,- vomito, as it is called,- which scourges these coasts, continues from the spring to the autumnal equinox, when it is checked by the cold winds that descend from Hudson's Bay. These winds in the winter season frequently freshen into tempests, and, sweeping down the Atlantic coast and the winding Gulf of Mexico, burst with the fury of a hurricane on its unprotected shores, and on the neighbouring West India islands. Such are the mighty spells with which Nature has surrounded this land of enchantment, as if to guard the golden treasures locked up within its bosom. The genius and enterprise of man have proved more potent than her spells.
After passing some twenty leagues across this burning region, the traveller finds himself rising into a purer atmosphere. His limbs recover their elasticity. He breathes more freely, for his senses are not now oppressed by the sultry heats and intoxicating perfumes of the valley. The aspect of nature, too, has changed, and his eye no longer revels among the gay variety of colours with which the landscape was painted there. The vanilla, the indigo, and the flowering cocoa-groves disappear as he advances. The sugar-cane and the glossy-leaved banana still accompany him; and, when he has ascended about four thousand feet, he sees in the unchanging verdure, and the rich foliage of the liquid-amber tree, that he has reached the height where clouds and mists settle, in their passage from the Mexican Gulf. This is the region of perpetual humidity; but he welcomes it with pleasure, as announcing his escape from the influence of the deadly vomito. He has entered the tierra templada, or temperate region, whose character resembles that of the temperate zone of the globe. The features of the scenery become grand, and even terrible. His road sweeps along the base of mighty mountains, once gleaming with volcanic fires, and still resplendent in their mantles of snow, which serve as beacons to the mariner, for many a league at sea. All around he beholds traces of their ancient combustion, as his road passes along vast tracts of lava, bristling in the innumerable fantastic forms into which the fiery torrent has been thrown by the obstacles in its career. Perhaps, at the same moment, as he casts his eye down some steep slope, or almost unfathomable ravine, on the margin of the road, he sees their depths glowing with the rich blooms and enamelled vegetation of the tropics. Such are the singular contrasts presented, at the same time, to the senses, in this picturesque region!
Still pressing upwards, the traveller mounts into other climates favourable to other kinds of cultivation. The yellow maize, or Indian corn, as we usually call it, has continued to follow him up from the lowest level; but he now first sees fields of wheat, and the other European grains, brought into the country by the conquerors. Mingled with them he views the plantations of the aloe or maguey (agave Americana), applied to such various and important uses by the Aztecs. The oaks now acquire a sturdier growth, and the dark forests of pine announce that he has entered the tierra fria, or cold region, the third and last of the great natural terraces into which the country is divided. When he has climbed to the height of between seven and eight thousand feet, the weary traveller sets his foot on the summit of the Cordillera of the Andes,- the colossal range that, after traversing South America and the Isthmus of Darien, spreads out, as it enters Mexico, into that vast sheet of tableland which maintains an elevation of more than six thousand feet, for the distance of nearly two hundred leagues, until it gradually declines in the higher latitudes of the north.
Across this mountain rampart a chain of volcanic hills stretches, in a westerly direction, of still more stupendous dimensions, forming, indeed, some of the highest land on the globe. Their peaks, entering the limits of perpetual snow, diffuse a grateful coolness over the elevated plateaus below; for these last, though termed "cold," enjoy a climate, the mean temperature of which is not lower than that of the central parts of Italy. The air is exceedingly dry; the soil, though naturally good, is rarely clothed with the luxuriant vegetation of the lower regions. It frequently, indeed, has a parched and barren aspect, owing partly to the greater evaporation which takes place on these lofty plains, through the diminished pressure of the atmosphere; and partly, no doubt, to the want of trees to shelter the soil from the fierce influence of the summer sun. In the time of the Aztecs, the tableland was thickly covered with larch, oak, cypress, and other forest trees, the extraordinary dimensions of some of which, remaining to the present day, show that the curse of barrenness in later times is chargeable more on man than on nature. Indeed the early Spaniards made as indiscriminate war on the forests as did our Puritan ancestors, though with much less reason. After once conquering the country, they had no lurking ambush to fear from the submissive semi-civilised Indian, and were not, like our forefathers, obliged to keep watch and ward for a century. This spoliation of the ground, however, is said to have been pleasing to their imaginations, as it reminded them of the plains of their own Castile,- the tableland of Europe; where the nakedness of the landscape forms the burden of every traveller's lament, who visits that country.
Midway across the continent, somewhat nearer the Pacific than the Atlantic ocean, at an elevation of nearly seven thousand five hundred feet, is the celebrated Valley of Mexico. It is of an oval form, about sixty-seven leagues in circumference, and is encompassed by a towering rampart of porphyritic rock, which nature seems to have provided, though ineffectually, to protect it from invasion.
The soil, once carpeted with a beautiful verdure and thickly sprinkled with stately trees, is often bare, and, in many places, white with the incrustation of salts, caused by the draining of the waters. Five lakes are spread over the Valley, occupying one tenth of its surface. On the opposite borders of the largest of these basins, much shrunk in its dimensions since the days of the Aztecs, stood the cities of Mexico and Tezcuco, the capitals of the two most potent and flourishing states of Anahuac, whose history, with that of the mysterious races that preceded them in the country, exhibits some of the nearest approaches to civilisation to be met with anciently on the North American continent.
Of these races the most conspicuous were the Toltecs. Advancing from a northerly direction, but from what region is uncertain, they entered the territory of Anahuac, probably before the close of the seventh century. Of course, little can be gleaned, with certainty, respecting a people whose written records have perished, and who are known to us only through the traditionary legends of the nations that succeeded them. By the general agreement of these, however, the Toltecs were well instructed in agriculture, and many of the most useful mechanic arts; were nice workers of metals; invented the complex arrangement of time adopted by the Aztecs; and, in short, were the true fountains of the civilisation which distinguished this part of the continent in later times. They established their capital at Tula, north of the Mexican Valley, and the remains of extensive buildings were to be discerned there at the time of the Conquest. The noble ruins of religious and other edifices, still to be seen in various parts of New Spain, are referred to this people, whose name, Toltec, has passed into a synonym for architect. Their shadowy history reminds us of those primitive races, who preceded the ancient Egyptians in the march of civilisation; fragments of whose monuments, as they are seen at this day, incorporated with the buildings of the Egyptians themselves, give to these latter the appearance of almost modern constructions.
After a period of four centuries, the Toltecs, who had extended their sway over the remotest borders of Anahuac, having been greatly reduced, it is said, by famine, pestilence, and unsuccessful wars, disappeared from the land as silently and mysteriously as they had entered it. A few of them still lingered behind, but much the greater number, probably, spread over the region of Central America and the neighbouring isles; and the traveller now speculates on the majestic ruins of Mitla and Palenque as possibly the work of this extraordinary people.
After the lapse of another hundred years, a numerous and rude tribe, called the Chichemecs, entered the deserted country from the regions of the far North-west. They were speedily followed by other races, of higher civilisation, perhaps of the same family with the Toltecs, whose language they appear to have spoken. The most noted of these were the Aztecs, or Mexicans, and the Acolhuans. The latter, better known in later times by the name of Tezcucans, from their capital, Tezcuco, on the eastern border of the Mexican lake, were peculiarly fitted, by their comparatively mild religion and manners, for receiving the tincture of civilisation which could be derived from the few Toltecs that still remained in the country. This, in their turn, they communicated to the barbarous Chichemees, a large portion of whom became amalgamated with the new settlers as one nation.
Availing themselves of the strength derived, not only from the increase of numbers, but from their own superior refinement, the Acolhuans gradually stretched their empire over the ruder tribes in the north; while their capital was filled with a numerous population, busily employed in many of the more useful and even elegant arts of a civilised community. In this palmy state, they were suddenly assaulted by a warlike neighbour, the Tepanecs, their own kindred, and inhabitants of the same valley as themselves. Their provinces were overrun, their armies beaten, their king assassinated, and the flourishing city of Tezcuco became the prize of the victor. From this abject condition the uncommon abilities of the young prince Nezahualcoyotl, the rightful heir to the crown, backed by the efficient aid of his Mexican allies, at length redeemed the state, and opened to it a new career of prosperity, even more brilliant than the former.
The Mexicans, with whom our history is principally concerned, came also, as we have seen, from the remote regions of the north,- the populous hive of nations in the New World, as it has been in the Old. They arrived on the borders of Anahuac towards the beginning of the thirteenth century, some time after the occupation of the land by the kindred races. For a long time they did not establish themselves in any permanent residence; but continued shifting their quarters to different parts of the Mexican Valley, enduring all the casualties and hardships of a migratory life. On one occasion, they were enslaved by a more powerful tribe; but their ferocity soon made them formidable to their masters. After a series of wanderings and adventures, which need not shrink from comparison with the most extravagant legends of the heroic ages of antiquity, they at length halted on the south-western borders of the principal lake, in the year 1325. They there beheld, perched on the stem of a prickly pear, which shot out from the crevice of a rock that was washed by the waves, a royal eagle of extraordinary size and beauty, with a serpent in his talons, and his broad wings open to the rising sun. They hailed the auspicious omen, announced by an oracle as indicating the site of their future city, and laid its foundations by sinking piles into the shallows; for the low marshes were half buried under water. On these they erected their light fabrics of reeds and rushes; and sought a precarious subsistence from fishing, and from the wild fowl which frequented the waters, as well as from the cultivation of such simple vegetables as they could raise on their floating gardens. The place was called Tenochtitlan, though only known to Europeans by its other name of Mexico, derived from their war-god, Mexitli. The legend of its foundation is still further commemorated by the device of the eagle and the cactus, which form the arms of the modern Mexican republic. Such were the humble beginnings of the Venice of the Western World.
The forlorn condition of the new settlers was made still worse by domestic feuds. A part of the citizens seceded from the main body, and formed a separate community on the neighbouring marshes. Thus divided, it was long before they could aspire to the acquisition of territory on the main land. They gradually increased, however, in numbers, and strengthened themselves yet more by various improvements in their polity and military discipline, while they established a reputation for courage as well as cruelty in war, which made their name terrible throughout the Valley. In the early part of the fifteenth century, nearly a hundred years from the foundation of the city, an event took place which created an entire revolution in the circumstances, and, to some extent, in the character of the Aztecs. This was the subversion of the Tezcucan monarchy by the Tepanecs, already noticed. When the oppressive conduct of the victors had at length aroused a spirit of resistance, its prince, Nezahualcoyotl, succeeded, after incredible perils and escapes, in mustering such a force, as, with the aid of the Mexicans, placed him on a level with his enemies. In two successive battles these were defeated with great slaughter, their chief slain, and their territory, by one of those sudden reverses which characterise the wars of petty states, passed into the hands of the conquerors. It was awarded to Mexico, in return for its important services.
Then was formed that remarkable league, which, indeed, has no parallel in history. It was agreed between the states of Mexico, Tezcuco, and the neighbouring little kingdom of Tlacopan, that they should mutually support each other in their wars, offensive and defensive, and that, in the distribution of the spoil, one fifth should be assigned to Tlacopan, and the remainder be divided, in what proportions is uncertain, between the other powers. The Tezcucan writers claim an equal share for their nation with the Aztecs. But this does not seem to be warranted by the immense increase of territory subsequently appropriated by the latter. And we may account for any advantage conceded to them by the treaty, on the supposition, that however inferior they may have been originally, they were, at the time of making it, in a more prosperous condition than their allies, broken and dispirited by long oppression. What is more extraordinary than the treaty itself, however, is the fidelity with which it was maintained. During a century of uninterrupted warfare that ensued, no instance occurred where the parties quarrelled over the division of the spoil, which so often makes shipwreck of similar confederacies among civilised states.
The allies for some time found sufficient occupation for their arms in their own valley; but they soon overleaped its rocky ramparts, and by the middle of the fifteenth century, under the first Montezuma, had spread down the sides of the tableland to the borders of the Gulf of Mexico. Tenochtitlan, the Aztec capital, gave evidence of the public prosperity. Its frail tenements were supplanted by solid structures of stone and lime. Its population rapidly increased. Its old feuds were healed. The citizens who had seceded were again brought under a common government with the body, and the quarter they occupied was permanently connected with the parent city; the dimensions of which, covering the same ground, were much larger than those of the modern capital.
Fortunately, the throne was filled by a succession of able princes, who knew how to profit by their enlarged resources and by the martial enthusiasm of the nation. Year after year saw them return, loaded with the spoils of conquered cities, and with throngs of devoted captives, to their capital. No state was able long to resist the accumulated strength of the confederates. At the beginning of the sixteenth century, just before the arrival of the Spaniard, the Aztec dominion reached across the continent from the Atlantic to the Pacific; and, under the bold and bloody Ahuitzotl, its arms had been carried far over the limits already noticed as defining its permanent territory, into the farthest corners of Guatemala and Nicaragua. This extent of empire, however limited in comparison with that of many other states, is truly wonderful, considering it as the acquisition of a people whose whole population and resources had so recently been comprised within the walls of their own petty city; and considering, moreover, that the conquered territory was thickly settled by various races, bred to arms like the Mexicans, and little inferior to them in social organisation. The history of the Aztecs suggests some strong points of resemblance to that of the ancient Romans, not only in their military successes, but in the policy which led to them.
1. Extensive indeed, if we may trust Archbishop Lorenzana, who tells us, "It is doubtful if the country of New Spain does not border on Tartary and Greenland;--by the way of California, on the former, and by New Mexico, on the latter!" Historia de Nueva España, (México, 1770,) p. 38, nota.
2. I have conformed to the limits fixed by Clavigero. He has, probably, examined the subject with more thoroughness and fidelity than most of his countrymen, who differ from him, and who assign a more liberal extent to the monarchy. (See his Storia Antica del Messico, (Cesena, 1780,) dissert. 7.) The Abbé, however, has not informed his readers on what frail foundations his conclusions rest. The extent of the Aztec empoire is to be gathered from the writings of historians since the arrival of the Spaniards, and from the picture-rolls of tribute paid by the conquered cities; both sources extremely vague and defective. See the MSS. of the Mendoza collection, in Lord Kingsborough's magnificent publication (Antiquities of Mexico, comprising Facsimiles of Ancient Paintings and Hieroglyphics, together with the Monuments of New Spain. London, 1830). The difficulty of the inquiry is much increased by the fact of the conquests having been made, as will be seen hereafter, by the united arms of three powers, so that it is not always easy to tell to which party they eventually belonged. The affair is involved in so much uncertainty, that Clavigero, notwithstanding the positive assertions in his text, has not ventured, in his map, to define the precise limits of the empire, either towards the north, where it mingles with the Tezcucan empire, or towards the south, where, indeed, he has fallen into the egregious blunder of asserting, that, while the Mexican territory reached to the fourteenth degree, it did not include any portion of Guatemala. (See tom. I. p. 29, and tom. IV. dissert. 7.) The tezcucan chronicler, Ixtlilxochitl, puts in a sturdy claim for the paramount empire of his own nation. Historia Chichemeca, MS., cap. 39, 53, et alibi.
3. Eighteen to twenty thousand, according to Humboldt, who considers the mexican territory ot have been the same with that occupied by the modern intendancies of Mexico, Puebla, Vera Cruz, Oaxaca, and Valladolid. (Essai Politique sur le Royaume de Nouvell Espagne, (Paris, 1825,) tom. I. p. 196.) This last, however, was all, or nearly all, included in the rival kingdom of Mechoacan, as he himself more correctly states in another part of his work. Comp. tom. II. p. 164.
4. The traveller, who enters the country across the dreary sand-hills of Vera Cruz, will hardly recognise the truth of the above description. He must look for it in other parts of the tierra caliente. Of recent tourists, no one has given a more gorgeous picture of the impressions made on his senses by these sunny regions than Latrobe, who came on shore at Tampico; (Rambler in Mexico, (New York, 1836) chap. 1;) a traveller, it may be added, whose descriptions of a man and nature, in our own country, where we can judge, are distinguished by a sobriety and fairness that entitle him to confidence in his delineation of other countries.
5. This long extent of country varies in elevation from 5570 to 8856 feet,--equal to the height of the passes of Mount Cenis, or the Great St. Bernard. The table-land stretches still three hundred leagues further, before it declines to a level of 2624 feet. Humboldt, Essai Politique, tom. I. pp. 157, 255.
6. About 62° Fahrenheit, or 17° Réaumur. (Humboldt, Essai Politique, tom I. p. 273.) The more elevated plateaus of the table-land, as the Valley of Toluca, about 8500 feet above the sea, have a stern climate, in which the thermometer, during a great part of the day, rarely rises beyond 45° F. Idem, (loc. cit.,) and Malte-Brun, (Universal Geography, Eng. Trans., book 83,) who is, indeed, in this part of his work, but an echo of the former writer.
7. The elevation of the Castiles, according to the authority repeatedly cited, is about 350 toises, or 2100 feet above the ocean. (Humboldt's Dissertation, apud Laborde, Itinéraire Descriptif de l'Espagne, (Paris, 1827,) tom. I. p. 5.) It is rare to find plains in Europe of so great a height.
8. Archbishop Lorenzana estimates the circuit of the Valley at ninety leagues, correcting at the same time the statement of Corté, which puts it at seventy, very near the truth, as appears from the result of M. de Humboldt's measurement, cited in the text. Its length is about eighteen leagues, by twelve and a half in breadth. (Humboldt, Essai Politique, tom. II. p. 29.--Lorenzana, Hist. de Neuva España, p. 101.) Humboldt's map of the Valley of Mexico forms the third in his "Atlas Geographique et Physique," and, like all the others in the collection, will be found of inestimable value to the traveller, the geologist, and the historian.
9. Humboldt, Essai Politique, tom. II. pp. 29, 44-49.--Malte Brun, book 85. This latter geographer assigns only 6700 feet for the level of the Valley, contradicting himself, (Comp. book 83,) or rather, Humboldt, to whose pages he helps himself, plenis manibus, somewhat too liberally, indeed, for the scanty references at the bottom of his page.
10. Torquemada accounts, in part, for this diminution, by supposing, that, as God permitted the waters, which once covered the whole earth, to subside, after mankind had been nearly exterminated for their iniquities, so he allowed the waters of the Mexican lake to subside in token of good-will and reconciliation, after the idolatrous races of the land had been destroyed by the Spaniards! (Monarchía Indiana, (Madrid, 1723,) tom. I. p. 309.) Quite as probable, if not as orthodox an explanation, may be found in the active evaporation of these upper regions, and in the fact of an immense drain having been constructed, during the lifetime of the good father, to reduce the waters of the principal lake, and protect the capital from inundation.
11. Anahuac, according to Humboldt, comprehended only the country between the 14th and 21st degrees of N. latitude. (Essai Politique, tom. I. p. 197.) According to Clavigero, it included nearly all since known as New Spain. (Stor. del Messico, tom. I. p. 27.) Veytia uses it, also, as synonymous with New Spain. (Historia Antigua de Méjico, (Méjico, 1836,) tom. I. cap. 12) The first of these writers probably allows too little, as the latter do too much, for its boundaries. Ixtlilxochitl says it extended four hundred leagues south of the Otomie country. (Hist. Chichemeca, MS., cap. 73.) The word Anahuac signifies near the water. It was, probably, first applied to the country around the lakes in the Mexican Valley, and gradually extended to the remoter regions occupied by the Aztecs, and the other semicivilized races. Or, possibly, the name may have been intended, as Veytia suggests, (Hist. Antig., lib. 1, cap. 1) to denote the land between the waters of the Atlantic and Pacific.
12. Clavigero talks of Boturini's having written "on the faith of the Toltec historians." (Stor. del Messico, tom. I. p. 128.) But that scholar does not pretend to have ever met with a Toltec manuscript, himself, and had heard of only one in the possession of Ixtlilxochitl. (See his Idea de una Nueva Historia de la América Septentrional, (Madrid, 1746,) p. 110.) The latter writer tells us, that his account of the Toltec and Chichemec races was "derived from interpretation," (probably, of the Texcucan paintings,) "and from the traditions of old men"; poor authority for events which had passed, centuries before. Indeed, he acknowledges that their narratives were so full of absurdity and falsehood, that he was obliged to reject nine-tenths of them. (See his Relaciones, MS., no. 5.) The cause of truth would not have suffered much, probably, if he had rejected nine-tenths of the remainder.
13. Ixtlilxochitl, Hist. Chich., MS., cap. 2.--Idem, Relaciones, MS., no. 2.--Sahagun, Historia General de las Cosas de Nueva España, (México, 1829,) lib. 10, cap. 29.--Veytia, Hist. Antig., lib. 1, cap. 27.
14. Sahagun, Hist. de Nueva España, lib. 10, cap. 29.
15. Idem, ubi supra.--Torquemada, Monarch. Ind., lib. 1, cap. 14.
16. Description de l'Egypte, (Paris, 1809,) Antiquités, tom. I. cap. 1. Veytia has traced the migrations of the Toltecs with sufficient industry, scarcely rewarded by the necessarily doubtful credit of the results. Hist. Antig., lib. 2, cap. 21-33.
17. Ixtlilxochitl, Hist. Chich. MS., cap. 73.
18. Veytia, Hist. Antig., lib. 1, cap. 33.--Ixtlilxochitl, Hist. Chich., MS., cap. 3.--Idem, Relaciones, MS., no. 4, 5.--Father Torquemada--perhaps misinterpreting the Texcucan hieroglyphics--has accounted for this mysterious disappearance of the Toltecs, by such fee-faw-fum stories of giants and demons, as show his appetite for the marvellous was fully equal to that of any of his calling. See his Monarch. Ind., lib. 1, cap. 14.
19. Texcuco signifies "place of detention"; as several of the tribes who successively occupied Anahuac were said to have halted some time at the spot. Ixtlilxochitl, Hist. Chich., MS., cap. 10.
20. The historian speaks, in one page, of the Chichemecs' burrowing in caves, or, at best, in cabins of straw;--and, in the next, talks gravely of their señoras, infantas, and caballeros! Ibid., cap. 9, et seq.--Veytia, Hist. Antig., lib. 2. cap. 1-10.--Camargo, Historia de Tlascala, MS.
21. Ixtlilxochitl, Hist. Chich. MS., cap. 9-20.--Veytia, Hist. Antig., lib. 2, cap. 29-54.
22. These were the Colhuans, not Acolhuans, with whom Humboldt, and most writers since, have confounded them. See his Essai Politique, tom. I. p. 414, II. p. 37.
23. Clavigero gives good reasons for preferring the etymology of Mexico above noticed, to various others. (See his Stor. del Messico, tom. I. p. 168, nota.) The name Tenochtitlan signifies tunal (a cactus) on a stone. Esplicacion de la Col. de Mendoza, apud Antiq. of Mexico, vol. IV.
24. "Datur hæc venia antiquitati," says Livy, "ut, miscendo humana divinis, primordia urbium augustiora faciat." Hist., Præf.--See, for the above paragraph, Col. de Mendoza, plate 1, apud Antiq. of Mexico, vol. I.,--Ixtlilxochitl, Hist. Chich., MS., cap. 10,--Toribio, Historia de los Indios. MS., Part 3, cap. 8,--Veytia, Hist. Antig., lib. 2, cap. 15.--Clavigero, after a laborious examination, assigns the following dates to some of the prominent events noticed in the text. No two authorities agree on them; and this is not strange, considering that Clavigero--the most inquisitive of all--does not always agree with himself. (Compare his dates for the coming of the Acolhuans; tom. I. p. 147, and tom. IV. dissert. 2.)--
See his dissert. 2, sec. 12. In the last date, the one of most importance, he is confimed by the learned Veytia, who differs from him in all the others. Hist. Antig., lib. 2, cap. 15.
25. The loyal Tezcucan chronicler slaims the supreme dignity for his own sovereign, if not the greatest share of the spoil, by this imperial compact. (Hist. Chich., cap. 32.) Torquemada, on the other hand, claims one half of all the conquered lands for Mexico. (Monarch. Ind., lib. 2, cap. 40.) All agree in assigning only one fifth to Tlacopan; and Veytia (Hist. Antig., lib. 3, cap. 3) and Zurita (Rapport sur lex Différentes Classes de Chefs de la Nouvelle Espagne, trad. de Ternaux, (Paris, 1840,) p. 11), both very competent crtitcs, acquiesce in an equal division between the two principal states in the confederacy. An ode, still extant, of Nezahualcoyotl, in its Castilian version, bears testimony to the singular union of the three powers.
26. See the plans of the ancient and modern capital, in Bullocks' "Mexico," first edition. The original of the ancient map was obtained by that traveller from the collection of the unfortunate Boturini; if, as seems probable, it is the one indicated on page 13 of his Catalogue, I find no warrant for Mr. Bullock's statement, that it was the same prepared for Cortés by the order of Montezuma.
27. Clavigero, Stor. del Messico, tom. I. lib. 2.--Torquemada, Monarch. Ind., tom. I. lib. 2.--Boturini, Idea, p. 146.--Col. of Mendoza, Part 1, and Codex Telleriano-Remensis, apud Antiq. of Mexico, vols. I., VI. |
Enemies: A History of the FBI
Random House. 537 pp. $30
By Susan Rosenfeld
In the Author’s Note prefacing his book on the FBI, Tim Weiner describes Enemies as the “history of the Federal Bureau of Investigation as a secret intelligence service,” its major mission, according to Weiner, for most of the past hundred years. The book chronicles the “tug-of-war between national security and civil liberties”—except that as Weiner portrays the FBI, with rare exceptions, there is no tug-of-war. “Security” far outweighs civil liberties and the Constitution.
This book is not an objective study of FBI history. Instead it selects examples that bolster the contention that the FBI put its wars against anarchists, Communists, the New Left, and foreign and domestic terrorists ahead of any consideration for the Bill of Rights. Weiner concedes that proponents from all these groups actually committed acts of espionage or violence. But for the most part, he features perpetrators who were never punished.
Weiner also oversells the role that surveillance played in J. Edgar Hoover’s FBI and beyond. As former foreign counterintelligence (FCI) agent Robert Lamphere noted in The FBI-KGB War, “only a small fraction of the New York field office [in the 1940s]—fifty or sixty men out of a thousand—was concerned with Soviet espionage and few agents outside the squad really knew or cared much about Soviet spies.” Add to that, foreign counterintelligence work was secret and could go on for years without resulting in any arrests or glory for its agents. That discouraged them from pursuing careers in FCI. By the post-Hoover era, foreign counterintelligence had become a backwater where one could place agents with the least ability such as Richard Miller, the first FBI agent to be accused and convicted of espionage.
At the same time, Weiner either minimizes Bureau successes or turns them into reasons for criticism. Prior to the bombing of Pearl Harbor, for example, the FBI had identified potential Japanese, German, and Italian spies and saboteurs and secured their arrest. Hoover opposed the 1942 internment of West Coast Japanese because in Hoover’s mind, everyone who posed a danger had already been detained. Instead, Weiner chose to emphasize whatever illegal techniques the FBI used to identify some of these enemies.
Weiner faults the Cold War FBI for not arresting more Russian spies. However, sources opened in the 1990s reveal that the Soviets had to change tactics and even recalled some spy handlers back home when FBI surveillance compromised their ability to contact their assets. As current FBI historian John Fox has noted, “Espionage is a difficult crime to prove, and prosecution for espionage, therefore, is not the standard by which to judge the success of a counterintelligence program.”
During Hoover’s tenure as FBI director (1924-1972), the Bureau’s domestic surveillance programs as well as FCI often fell into a gray area not specifically addressed by the courts, and the attorney general’s formal guidelines regarding surveillance did not exist until 1976. The FBI also liberally interpreted decisions that might have compromised its surveillance programs. Section 605 of the Federal Communications Act of 1934, for example, did not permit the wiretapping and divulging of communications without the sender’s permission. Hoover decided that wiretap surveillance remained legal as long as its contents were not “divulged.” Even after the Keith case in 1972 specifically outlawed wiretaps on Americans, the courts refused to include foreign surveillance: “The instant case requires no judgment on the scope of the president's surveillance power with respect to the activities of foreign powers, within or without this country.”
After 1978, the Foreign Intelligence Surveillance Court (FISC) had to approve electronic surveillance. Weiner, like many Bureau critics, pointed out that the FISC rarely turned down a Bureau request. Actually, the briefs requesting surveillance usually went through many drafts until the FBI developed a case that the court would accept.
If one looks only at Hoover and the FBI, naturally, it will seem as if they are responsible for everything. Weiner usually avoids this by emphasizing that presidents urged the Bureau to conduct surveillance by any means necessary. Nevertheless, Weiner alleges that “Hoover was creating the political culture of the Cold War in the United States.” But a number of other individuals and organizations feared American Communists and Soviet spies as much as Hoover. According to Richard Gid Powers, “American anticommunism [after World War II] moved from the margins to the center of American politics, as Joseph Stalin . . . extended his power across half of Europe.” The Soviet Union’s aspiration of spreading Communism throughout the world was a fact. Even the liberals of Americans for Democratic Action (ADA), as well as many other religious and veterans groups, perceived Communism as a potential threat to American democracy. They did not need J. Edgar Hoover and the FBI to determine that Communism and Russian totalitarianism should be contained. Powers, incidentally, is himself a historian of the FBI, and author of what I consider to be the best biography of J. Edgar Hoover, Secrecy and Power (1986). If Weiner was aware of Powers’s books, he did not cite any of them.
The anti-war, feminist, and black power movements of the 1960s and 1970s, dubbed the “New Left” in the FBI, became the focus of both FBI intelligence and counterintelligence programs. In these investigations, the FBI infiltrated suspect organizations using informants and undercover agents, collected documents, and observed demonstrations. The military, Secret Service, Central Intelligence Agency (CIA), and local law enforcement agencies also used these methods, although the FBI was almost always the agency blamed for such intrusions. Interspersing the “spying” and counterintelligence programs against the New Left with unrelated investigations, as Weiner does, makes them lose their impact. Yet their revelation in the post-Watergate years, along with that of Hoover’s vendetta against Martin Luther King, Jr., forever damaged the reputation of Hoover and the Bureau, leading to the kind of distortion of FBI history exemplified by Enemies. Weiner, like many commentators on surveillance of the New Left, neglects the bombings, property destruction, and threats of violent revolution coming from members of some of these organizations.
Regarding King, Weiner correctly identifies the assistant director in charge of the Domestic Intelligence Division, William Sullivan, father of the infamous Counterintelligence Program (COINTELPRO), as the author of the note to King that may have suggested that the civil rights leader commit suicide—or possibly just turn down the Nobel Prize. King did neither, of course, and Hoover is often erroneously credited with authoring the note. Weiner, however, distorts the Hoover-King situation elsewhere by implying that the FBI director meant King’s sex life when he declared the civil rights leader “the most notorious liar in the country.” Rather, according to one of Hoover's major advisors, Cartha D. “Deke” DeLoach, this remark referred to King’s failure to withdraw his inaccurate accusation that most FBI agents assigned to the South were themselves southerners and therefore prejudiced against Negroes.
The FBI also conducted one of its largest investigations ever to find King’s assassin. To identify and then track down James Earl Ray took painstaking work and detective skills. Yet Weiner makes it appear that Scotland Yard solved the crime after an inept Bureau let Ray escape. The investigative files are available on the FBI’s website. If Weiner read the files his characterization of the investigation represents a deliberate distortion. If he did not, it represents a serious research failure. This section is only one example in which Weiner does not cite what sources he used.
Hoover died on May 2, 1972. Approximately six weeks later, burglars linked to Nixon’s re-election committee broke into Democratic headquarters in the Watergate office complex in Washington, DC. Not only did Watergate lead to the resignation of a president, it precipitated Congressional and media investigations into America’s intelligence community. The resulting changes in the nation’s security apparatus made intelligence collection more difficult. Furthermore, the revelations about break-ins and wiretaps invoked fear that the FBI’s principal activity was spying for the sake of spying and for very little reason, a critique Weiner echoes. Practices like following aging Communists or placing informants in women’s consciousness-raising sessions opened the FBI to ridicule. Neither helped the United States battle foreign or domestic terrorism, nor to identify spies within US organizations. In addition, after the fall of the Soviet Union, then-FBI Director William S. Sessions downgraded counterespionage. According to Weiner, Hoover-era-type tactics would not be revived until after the 9/11 terrorist attacks. “For years to come,” he writes, “the FBI agents who hunted terrorists in America wandered in a legal wilderness, looking for signs to guide them through an uncharted land.”
For the next hundred pages or so, Weiner compiles Bureau attempts to vanquish terrorists while maintaining strict adherence to the law. Fifty-three agents, according to Weiner, were investigated by the Department of Justice for illegal breaches of privacy (no source given). The Bureau did not protect its agents, making counterterrorism an assignment to avoid. Of course, some agents diligently continued their investigations, such as the ultimately successful resolution of the 1976 assassination in Washington of Chilean exile Orlando Letelier. Weiner buries such successes among failures, for example, by featuring spies within the FBI like Robert Hanssen. Little of these discussions concerned civil liberties questions. In fact, few of the post-Hoover sections of the book deal with civil liberties and security. Weiner apparently had other stories to tell, especially if they put the FBI and its directors in a negative light.
Enemies is not a page-turner. Rather, as Weiner states in the Author’s Note, it is a compendium of “illegal arrests and detentions, break-ins, burglaries, wiretapping, and bugging,” usually justified by a presidential request endorsed tacitly or actually by the attorney general. Most of Weiner’s examples run a few sentences or paragraphs. He rarely tells a story in full, with the exception of the 1919 Red Raids, Watergate, and strangely, the FBI’s role in suppressing a revolt in the Dominican Republic in 1965, which rates an entire chapter.
For as experienced a journalist as Weiner, who has worked for the Philadelphia Inquirer and New York Times, to enliven his text with melodramatic descriptions of Hoover’s (and others) supposed feelings is surprising. One example: Hoover’s “rage at the president’s reluctance to fight a full-bore war on communism grew ferocious.” Hoover’s suggesting that he wanted a “showdown” with the president (Truman), and asking members of Congress to “give him the power to protect America against ‘the threat of infiltrating foreign agents, ideologies, and military conquest,’” does not sound to me like a “ferocious” rage. Rather Hoover was convinced that his FBI could do a better job than any other government agency in protecting the nation against foreign and domestic enemies.
Weiner uses an impressive array of sources in the 60 pages of notes, almost all of them primary documents, including some declassified as recently as August and November 2011. He also incorporates oral histories and interviews from the Society of Former Special Agents of the FBI’s Oral History Heritage Program, many of which are available to the public on the National Law Enforcement Officers Memorial Museum website. (I am a consultant to the project and at least one of my interviews is cited.) In addition, Weiner graciously points to several secondary works that pulled together new information, notably those by Raymond J. Batvinis on the Special Intelligence Service, and the work of Katherine A. S. Sibley, John Earl Haynes, Harvey Klehr, and Alexander Vassilev on Soviet espionage.
The scope of his sources and the many pages of citations, nevertheless, leave a number of problems that seriously diminish the credibility of this book. In most cases, Weiner cites only direct quotations. As a former official historian of the FBI, I want to know the origin of controversial vignettes or of information that is new to me. Often Weiner does not give an adequate source. Sometimes whole paragraphs go by without any attribution.
Still another frustration for me was in identifying exactly which sources Weiner used when he does cite them. Some notes are quite specific and therefore quite helpful. In his discussion of the FBI’s early years (the records of which are on microfilm), he gives both the reel number and the case number. Weiner also includes links to some specific online documents. But in citing the many FBI Freedom of Information Act releases, he uses only the abbreviation “FBI/FOIA.” The file names are not cited, making it virtually impossible to locate the document to check the context of a quote. Moreover, he does not indicate whether the source is from records he recently acquired or among those available on the FBI’s website.
Weiner also appears to accept at face value oral history interviews that he cites. In the section on Watergate, he includes former Special Agent Daniel F. Bledsoe’s verbatim recollection of a conversation that allegedly took place the day of the break-in. Bledsoe was the FBI headquarters supervisor on duty that weekend. In the first place, Weiner only partially reconstructs and sometimes paraphrases Bledsoe’s rendition. Second, Bledsoe was recalling an event that took place almost 37 years earlier. Any good researcher should know that recollections are often fallible, and will attempt to confirm any major allegations. Bledsoe describes emphatically and in detail that presidential aide John Ehrlichman ordered him to stop the FBI from investigating the Watergate break-in and threatened to have Bledsoe fired when he refused. That certainly constitutes a major allegation!
Why a high-level White House official would give such an order to an FBI supervisor and not to its highest executives should have raised a question. Apparently it did not. Had Weiner wanted to probe further, he could have found Bledsoe’s handwritten notes from 17 June 1972 in the FBI’s Watergate files available online. The closest the notes come to any instruction from the White House is in a notation relaying a message from the FBI’s Washington Field Office. Chief of staff H. R. Haldeman’s aide, Alexander Butterfield, had informed Ehrlichman of the break-in, and in turn, an order had been issued that “No abnormal pressure [be] put on CIA per . . . White House.” Did Bledsoe’s memory turn this third-hand mention of Ehrlichman into an anecdote that put Bledsoe in the middle of one of the crimes of the century? Did such a call from Ehrlichman or another official occur at another time to Bledsoe or another FBI person? If it did, it has not found its way into any primary or secondary source that I am aware of, although the former agent who interviewed Bledsoe noted, “I’ve heard that a couple of times.” Regardless, the episode presented significant new information. It should have elicited a thorough inquiry by Weiner.
Enemies does not present a 100 percent negative picture of the FBI even under Hoover. For example, Weiner summarizes Hoover’s accomplishments from the 1920s and 1930s in one paragraph: “He fired crooks and incompetents . . . instituted uniform crime reports, built a training academy, and assembled a national fingerprint file.” Despite Hoover’s centrality to this book, Weiner does not make the mistake of equating the entire FBI with its director. Those who do, cannot understand, as Weiner does, that some FBI agents aggressively pursued criminal and terrorist investigations against the Klan and “White Hate” groups at the same time other agents were investigating various black civil rights groups for communist infiltration.
As Hoover neared the end of his reign, Weiner credits him with refusing to sign the Nixon administration’s plan to remove many of the legal restrictions on the various intelligence services. Hoover’s refusal convinced the heads of the other intelligence agencies to decline as well. As a result, the White House developed its own agents, aka the “plumbers,” which in turn, led to Watergate. Weiner also notes Judge William H. Webster’s efforts as director to “do the work the American people expected in the way that the Constitution demanded.” Current FBI Director Robert S. Mueller is a hero to Weiner because of his refusal to agree to President George W. Bush’s request for unlimited domestic surveillance. At the book’s end, Weiner trusts Mueller to achieve a proper balance between freedom and security in the future.
Weiner does not make a specific effort in his narrative to relate the past to the present. However, readers interested in his topic cannot help but make the connection. To many readers, this book will remind Americans they must remain vigilant if they are to maintain the civil liberties. I do not dispute that important goal. But by demonizing the FBI, authors such as Weiner, in my opinion, go too far. Even J. Edgar Hoover made an effort to justify the “illegalities” his agents performed in the name of national security. Today Congress and the media criticize the Bureau for not “connecting the dots,” and at the same time civil libertarians bring up the FBI’s history of over-reaching surveillance. Weiner purported to portray the difficulties in balancing this conflict. Instead, he wrote a catalog of the Bureau’s alleged intelligence misdeeds, and a critique of the post-Watergate FBI that had little to do with perceived constitutional breaches.
Susan Rosenfeld served as the FBI’s first official historian from 1984 to 1992. Since 2002, she has been the principal consultant to the Society of Former Special Agents Oral History Heritage Program. Earlier in her career she was an archivist at the National Archives and Records Administration, specializing in Department of Justice documents, including the records of the FBI and the Watergate Special Prosecution Force. She received her PhD in history from Georgetown University, and has also taught US constitutional history. Currently she is a historian with the US Air Force, Air National Guard.
This review represents the personal opinion of its author, and not those of her former employers, including the Federal Bureau of Investigation, or her present employer, the US Air Force.
Tim Weiner, Enemies: A History of the FBI (New York: Random House, 2012), xv.
Robert J. Lamphere and Tom Shachtman, The FBI-KGB War: A Special Agent’s Story (Macon, GA: Mercer University Press, 1995), 20; Weiner, Enemies, 356-357.
John Fox, “What the Spiders Did: US and Soviet Counterintelligence Before the Cold War,” Journal of Cold War Studies, Vol. 11, No. 3 ( Summer 2009), 220; note, 222.
United States v. United States District Court, 407 U.S. 297. The quote is from Section II.
Royce Lamberth, “FISA Court Judge Royce Lamberth Discusses Work of Court.” National Security Law Report Vol. 19, No. 2 (May 1997), 1-2, 4-5.
Weiner, Enemies, 146; Richard Gid Powers, Not Without Honor: The History of American Anticommunism, (New Haven, CT: Yale University Press, 1995), 191. Powers’s books on the FBI include G-Men: Hoover’s FBI in American Popular Culture (Carbondale, IL: Southern Illinois University Press, 1983); Secrecy and Power: The Life of J. Edgar Hoover (Free Press, 1987); and Broken: The Troubled Past and Uncertain Future of the FBI (New York: Free Press, 2004).
Weiner, Enemies, e.g., 270-273; 279-280.
Ibid., 250; Cartha D. “Deke” DeLoach, Hoover’s FBI: The Inside Story by Hoover’s Trusted Lieutenant (Washington, DC: Regnery Publishing, 1995), 204.
Weiner, Enemies, 274.
Ibid., 328.
Ibid., xvi.
Ibid., 146.
Raymond J. Batvinis, The Origins of FBI Counterintelligence (Lawrence, KS: University Press of Kansas, 2007); Katherine A. S. Sibley, Red Spies in America: Stolen Secrets and the Dawn of the Cold War (Lawrence, KS: University Press of Kansas, 2007); John Earl Haynes, Harvey Klehr, and Alexander Vassiliev, Spies: The Rise and Fall of the KGB in America (New Haven: Yale University Press, 2006).
Weiner, Enemies, 309-310; Society of Former Special Agents of the FBI, Inc., Interview of former Special Agent of the FBI Daniel F. Bledsoe (1955-1980), interviewed by Brian R. Hollstein, 19 August 2009.
Weiner, Enemies, 63, 248.
Ibid., 291, 296, 344, 448.
© 2012 by Susan Rosenfeld |
When he pitched a hitless inning in relief for the New York Highlanders (already beginning to be called the Yankees) on April 16, 1907, Roy Castleton became the first native of Utah to play major league baseball. Castleton, however, was much more than the answer to a trivia question. He was a promising left-handed pitcher with a minor league perfect game to his credit.
The story of the Castleton family is the story of Utah's settlement. Between 1860 and 1870, over 9,000 persons emigrated from England to Utah. One of those families was the Castletons. The family had lived in Lowestoft, Suffolk, England, since at least the mid-eighteenth century. James Joseph Castleton was born there in 1829, the eldest of ten children, and worked as a fisherman as well as a rope and twine maker. On January 2, 1854, he married Frances Sarah Brown, from the nearby village of Pulham. On October 4 of that year, their first son, Charles, was born.
During the 1850s, the Mormon Church wanted to expand and develop their Utah colony. Much of their effort was concentrated on gaining converts in Europe, with a particular emphasis on the British Isles. At least one member of the Castleton family, Frances, was among those converts. The church established the Perpetual Emigrating Fund Company, which helped more than a third of church members emigrating from Europe to the American West.
In early 1863, James Castleton, his wife and four sons were among the immigrants. They left London on the Amazon around the beginning of June. They reached the Port of New York on July 20. The family traveled by train to Omaha, Nebraska, then secured teams of oxen and set out across the plains, reaching the Salt Lake Valley on October 4. Frances was pregnant at the time and rode while those able, including Charles, walked much of the way across the West.
James Castleton was hired by Brigham Young as a gardener, and was still at that job in 1870. He saved enough money to buy land at L Street and Seventh Avenue and open a store.
Charles Castleton worked at the store, and became a carpenter. In 1879 he married Mary Ann Luff, also an immigrant from England. Charles and Mary Ann would have seven children. On July 26, 1885, a son, Royal Eugene, was born.
Roy grew up in a middle class neighborhood of Salt Lake City's Fourth Ward. The city itself had a population of around 45,000 during his childhood. It was also becoming a more diversified place to live. About half the population did not belong to the Latter-day Saints Church, though Roy and his family were members. When Roy was ten, Utah became the 45th state.
Castleton was well educated by turn-of-the-century standards: He was a high school graduate who especially enjoyed mathematics. His other passion was baseball.
By the summer of 1904 he was using both talents. He joined his older brother Charles as a member of the state's best amateur team, the Cleveland Commission company team of Salt Lake City. According to the Ogden Standard Examiner the team included several former professional players and was the undisputed amateur champion of the state. Not yet 18, Roy more than held his own in a series of games with an independent professional team from Ogden. On July 3, he lost a 6-5 fourteen-inning game. The Standard Examiner said Castleton "handled the horsehide in good style, and finished with a strong wing at the end of the fourteenth inning." Before a thousand fans, the largest crowd of the season in Ogden, he won the next day by the same 6-5 score, with his endurance noted. He went 3-4 at the plate that day.
At some point late in the 1904 season, Roy signed with his hometown Salt Lake City team in the Class B Pacific National League. He probably played at least briefly with the team in the closing weeks of the season, and was included on the team's reserve list that fall.
During the winter of 1904-05 Frank Gimlin, manager of the Cleveland Commission team, was hired by Ogden as the team entered the Pacific National League.
Gimlin recruited some of his former players, including acquiring Castleton from Salt Lake City. Only nineteen, Castleton seemed overmatched by the more experienced competition, but he made the team and pitched well at times. His best outing was a sixteen-inning loss at Salt Lake City.
After the league collapsed in June, the Ogden team stayed together as an independent team, and Roy also earned extra money pitching for a semipro team in Blackfoot, Idaho. In Blackfoot, he pitched under an assumed name, taking the name of Sweeney, the pitcher he'd replaced. The Los Angeles Times later told the story of Castleton's last appearance as "Sweeney." "Blackfoot fans ran a special excursion [train] down to Idaho Falls. Unfortunately, a traveling man from Salt Lake happened to be among the excursionists. He recognized Castleton. Just before the game started, a prominent citizen of Idaho Falls, with a Gatling gun in his pocket, stepped to the plate and announced that Blackfoot had a ringer in the box and declared all bets off. Neither side scored for three innings, and a lot of new bets were made in the meantime. Blackfoot won in a driving finish, and Castleton at once became the center of interest. Roy started for the depot on [sic] a gallop. He hastened in every respect of the word. Almost anyone would be prone to hasten with a large portion of the population of Idaho Falls and surrounding country in his wake and considerably peeved. After great difficulty, Castleton was hoisted on board a train for Blackfoot."
When not pitching or being chased by angry fans, Castleton worked as a clerk and bookkeeper for one of the railroad offices in Salt Lake City. In the spring of 1906, he took one of those trains east.
Youngstown, Ohio, had been a hotbed of professional baseball for several years. After a few seasons of independent ball, the city along with many others entered organized ball in July of 1905 as part of the Class C Ohio Pennsylvania League (O-P). Despite the jumble of teams and inconsistent scheduling, Youngstown had a strong team and was ruled league champion. By 1906, the league had shrunk to a more manageable eight teams, six of them in Ohio.
Although the league was new, O-P teams had established rosters largely consisting of veteran players. Most of the Youngstown players that spring were at least four or five years older than Castleton, the youngest player on the team. At a listed height of 5-11 1/2 and 150 pounds, he was one of the team's tallest players but looked like a man who'd missed a meal or two.
Youngstown manager Marty Hogan, a former outfielder, was an outstanding judge of pitching talent. Later he'd be responsible for signing Stan Coveleski and Sam Jones to their first minor league contracts, but in 1906 Castleton was his find.
It took just one exhibition game to impress the Youngstown Telegram: "Roy Castleton has a jump ball that has it on anything else seen in this neck of the woods. When he has the ball working right a batter has trouble in sending it to the outfield."
Castleton made his first regular season appearance in Youngstown's third game. He beat Newark 6-3 but had to be relieved after six innings. The Telegram said he was "plainly rattled and only the clever work in the outer garden held Youngstown safe." He gave up five hits, struck out two, and walked one. It was still a better debut than that made by another young southpaw a few days later. Richard Marquard, already nicknamed Rube, lasted just one third of an inning in a relief appearance for Lancaster, allowing five runs on seven hits.
Castleton won four of his first five decisions, two of those wins in relief, and the Telegram said, "with a little coaching by the veteran catcher, Lee Fohl, [Castleton] should develop into a real wonder. Roy has not been backing up the bases well in the games he has played, probably due to a little nervousness, but he will soon overcome this. He has an abundance of smoke, good curves and excellent control." The Telegram noted those curves were "breaking much better in practice than in the games in which he has participated but each game shows an improvement." The article predicted that he would be the best southpaw in the league by the end of May. That prediction wasn't far off.
Even when he was struggling on the mound, Castleton was one of the league's better gate attractions. He told a Utah newspaper that he was called "the tall Mormon pitcher, or Brigham Young. While on the coaching lines the fans were continually asking how many wives he had at home, all of which he took good naturedly. Whenever Castleton was slated to pitch, it was used as an advertisement by the home teams and generally resulted in bringing out a large delegation of girl fans who wanted to see the man with many wives."
July, Castleton's strongest month, marked the first mention of major league interest in the young southpaw. The sports editor of the Youngstown Telegram received inquiries from three unnamed major league managers, and Hogan was also contacted about the 20-year-old pitcher.
On the field, Castleton was developing the consistency the Telegram had predicted. He won six of seven starts and also won a game in relief. His record stood at 15-8 the morning of August 17. That game would make Roy Castleton famous.
It was the second game of a key series between first place Youngstown and second place Akron. Hogan felt Akron was weak against lefthanders and decided to pitch Roy outside his normal spot in the rotation. The next day's Telegram reported, "In the most remarkable game ever seen in Youngstown, remarkable for the great pitching of Roy Castleton, Youngstown shut out Akron 4-0. Akron never had a lookin [sic], not at a run, nor a hit, nor even first base. Castleton was in anything but a generous mood and although the visitors were of the opinion that he might 'open up' and permit one or more of them to get to first base, the Mormon youth thought otherwise and not a single Tip Top made the acquaintance of Mert Whitney at the initial sack."
The Telegram also compared Castleton's perfect game to the one thrown by Cy Young the year before. "In that memorable contest Young did not allow a run nor a hit. He did not give a base on balls nor hit a batsman. Castleton not only did this Friday but he went old Cy a few better. Only four balls were batted out of the infield. The Mormon also compelled ten men to fan the atmosphere." A dropped third strike was the closest Akron came to a baserunner, but catcher Fohl threw the runner out easily.
Major league managers took notice of the perfect game. When Youngstown went to Mansfield for their next series, one of those in attendance was Clark Griffith, manager of the New York Highlanders. He offered $2,000 for Castleton with the young pitcher reporting to New York immediately. The team's owners turned down the offer, as Youngstown was in a tight pennant race.
Roy was considered one of the most gentlemanly players in what wasn't a gentlemen's league. Still on one occasion late in the season his temper got the best of him. The Telegram noted, "Roy Castleton, the Mormon youth, is not immune when it comes to saying cuss words. Of course Roy uses them in a modified form but just the same he uses them. Umpire Sam Wise's eyesight was very bad, and he sent three men down to first in succession in the first inning, a run being forced across the plate. Both Castleton and Fohl claim that over one half of the balls cut the plate almost in two, and were of regulation height. Sam couldn't see them that way, however, and Roy expressed his opinion of Wise in anything but a complimentary manner." That bases loaded walk was decisive as Castleton lost the game to Akron 4-3.
With a month left in the O-P season, Griffith got his man. New York drafted Castleton for the standard Class C price of $500. Because of his being drafted Youngstown retained Roy for the rest of the season. He finished by winning his last five decisions, finishing with a 22-12 record and striking out 156 batters in 278 innings.
Castleton reported to New York's spring training site of Atlanta, Georgia, on Saturday, March 9. Alexander MacKenzie, sports editor of the New York Mail, said of Roy's early workouts: "The most promising [of the rookie pitchers], to my notion, is the young lefthander, Castleton. He has better control than any lefthander seen in the East for many years, and had so much speed that Griff had to stop him a couple of times during the practice."
Sid Mercer of the Globe also wrote about the young southpaw: "Jack Kleinow was doing the receiving. Presently the ball began to produce loud and resounding thumps in his big mitt. He was faster than it really looked possible for him to be. Then he began hooking them over, and Griff's eyes opened wider, for Castleton was throwing one of the best curves which Griff had ever seen so cleverly controlled by a young southpaw." Kleinow said of that curve, "A right handed batter will fall for that ball every time, it breaks so quickly, though, that he can't dodge, and even if he doesn't intend to hit, he will throw up his stick to protect himself. The usual result is either a pop fly or an easy grounder, with the runner getting a bad start to first."
After splitting their first two regular season games, New York opened at home against the Philadelphia Athletics. Despite bad weather, the New York Times estimated the crowd at 10,000. Five pitchers were used between the two teams that day. Roy Castleton was the third of three pitchers to appear for New York. Entering the game in the top of the ninth, he worked a hitless inning.
After nearly a week of inactivity, Castleton was the starting pitcher in an exhibition contest with Newark of the Eastern League. On a cold afternoon marked by occasional snow, he allowed thirteen hits in a 12-3 loss. Two days later he was optioned to Atlanta for more seasoning. Though he'd barely pitched in a month, Castleton's debut in the Southern Association was impressive. On April 29 he beat defending champion Birmingham 5-1. Roy allowed just four hits and struck out seven.
His third Atlanta appearance illustrated both his potential and his flaws as a pitcher. He was dominating early, striking out four batters in the first two innings. He finished with ten strikeouts in a thirteen-inning tie at Memphis, but struggled badly in the eighth inning. An error and a pair of walks that inning let Memphis score both their runs. He walked six in that game. He later appeared in two twelve-inning games, winning one and tying the other. Those were Atlanta's three longest games of the season.
Three losses and a brief period of inactivity with a sore arm slowed his progress through the last half of May. The Atlanta Constitution was still impressed with the young southpaw: "He was bumped on the road, but that was a natural occurrence; no one holds it against him, nor does anyone lessen the value placed on him when he was winning so regularly, and receiving the plaudits of the multitude. He promises to get more plaudits as well as bumps."
Not all the bumps he received were on the field. Shortly after joining Atlanta, an imaginative newspaper reporter claimed he was a Mormon with sixteen wives. A report a few years later in the Los Angeles Times said Castleton had been visited by "representatives from the various civic bodies, the W.C.T.U., and The Women's Foreign Missionary Society. If Roy actually had sixteen wives, it was deemed advisable to ride him out of town on a rail. Roy established his singleness to the satisfaction of all concerned and continued to remain in the circuit until called higher." Though certainly exaggerated, this incident showed some of the off-field distractions that Castleton had to face pitching in the South.
June and July of 1907 were successful months for Castleton. He won 10 of 13 decisions for Atlanta, developing into the league's best lefthander. One of his best midseason outings was on July 4. The Constitution said: "the southpaw had them on the on the blink throughout the nine rounds, pitching almost faultless ball the while. The small covey of hits [three] were scattered here and there, and in that condition were worthless." He was more effective on the road than at Atlanta's Ponce De Leon Park. He said the mica soil of the pitching mound had a negative effect on his control.
As well as Castleton pitched the first four months of the season, he saved his best work for the heat of the pennant race. On August 30, he shut out Little Rock on four hits. No runner reached third. Three days later he struck out nine Shreveport batters throwing a five-hit shutout, winning 5-0 in the second game of the Labor Day doubleheader. He threw his third straight shutout on September 6, blanking New Orleans by an identical 5-0 score. His final 1907 start for Atlanta was his fourth straight shutout. It was a classic pitching duel against Memphis' George Suggs. Suggs allowed five hits, striking out eight and walking two. Castleton was even better. He allowed just four hits with eight strikeouts and one walk. Appropriately the contest ended in a scoreless tie. His final record with Atlanta was 17-8, and he was generally considered the best pitcher on a staff that included Russ Ford and Bob Spade. When an all-time Atlanta team covering the years 1902-07 was picked the following spring, Castleton was among those chosen.
Griffith exercised his option on the young southpaw, and Castleton rejoined New York at the end of the Southern Association season. He made two starts during the last week of the American League season. His first start was in the opener of a doubleheader against the St. Louis Browns on September 28. The game was a classic deadball era game. Castleton retired the first twelve Browns he faced before Bobby Wallace opened the fifth with an infield single. St. Louis scored two in the sixth and added another in the seventh for a 3-1 win. Despite the loss, Castleton was very effective, allowing five hits, striking out two and walking one. The New York Times said of his first start, "The Southern League graduate soon installed himself as a favorite of the fans."
Castleton's second and, as it turned out, final American League start was on October 2 against Doc White and the Chicago White Sox. New York scored four runs in as many innings off White, and Castleton could be described as effectively wild that afternoon. Charles Dryden's account of the game in the Chicago Tribune commented on Castleton's control in humorous fashion. "After [Jiggs] Donahue doubled in the second Castleton filled the bags by hitting two athletes [Patsy Dougherty and Hub Hart] in vulnerable spots below the belt. By the time he had done all these heroic stunts two men were out, Dr. White mortified his colleagues by fanning the air which was cool enough."
Through the first five innings, Castleton allowed the defending World's Champs just three hits. Dryden's account described a Chicago three-run rally: "The Sox bubbled in the sixth. Donahue soaked a safety. [George] Davis fanned. [Charlie] Hickman [replacing the injured Dougherty] got hit on the excess baggage. [George] Rohe's smash to left bounded over the low bleacher fence for a home run." Dryden claimed the home run "got Castleton's goat," and after getting out of the inning with a one-run lead, "Castleton was so glad he was alive he at once turned over his portfolio to slow Joe Doyle." Doyle held Chicago scoreless, saving Roy's first major league win. Despite the three- run inning, Castleton still pitched well, allowing six hits, striking out three, walking one and of course hitting the three batters below the belt.
After the season ended, Castleton returned to Salt Lake City and his work in the railroad office. He later said, "I have to work. If my mind is not occupied, I get to thinking too much about things and consequently get blue. And I have to work at something that will keep me going hard. I never felt better than when I had whole sheets of figures to add in making trial balances. The more the merrier. And if I got hold of something with about fourteen columns a mile long each, and about five sheets of them, I was in my glory." He said he usually got off work early enough to keep in shape with an afternoon workout.
Despite Castleton's solid pitching, Griffith decided he needed a little more seasoning. In late February, Griffith optioned Castleton to Atlanta for the 1908 season. The option to Atlanta was meant as park rental for New York's spring training stay.
Castleton's 1907 salary was a controversial topic in the spring of 1908. The Memphis Commercial-Appeal claimed the Atlanta groundskeeper "was paid an enormous salary. This, of course, did not appear on the players' payroll, but it is charged was cut up with a star pitcher, who received his salary partly in check form, with the remainder in cash handed as a 'present' from the keeper of the grounds." Action was threatened against offending clubs by the league president, but nothing was proven.
An on-field spring highlight was a March 28 start against the Cubs. When Castleton left after five innings, the Crackers led the defending World Champions 3-2. After allowing a pair of runs in the first, Castleton was very effective. The Constitution reported, "He worked four more rounds and there were but two more hits, far apart as hits should be when the home team is not getting them."
Abnormally cool weather contributed to a slow 1908 start for the Atlanta southpaw. Exaggerating somewhat, the Constitution said of his first start against Nashville, "By the fourth round icicles began to form on his hurling wing." He pitched well in a loss that day but lasted just two innings in his second start. He didn't pitch for three weeks after that, probably due to injury.
When Roy returned to the mound at Little Rock on May 13, he had a new pitch that was helped by adverse weather conditions. The Constitution reported, "Heavy rains of the night preceding soaked through the crust at West End Park, and made the going bad, even for web footed and water-wagon stars of the diamond. The ball was slippery and spitters served by Castleton were not necessarily wetted by the Mormon's saliva." He allowed two hits and struck out ten that afternoon and had "local batters up Salt Lake all the way through."
He lost his next two starts to drop to 1-4. The second of those was an eleven-inning 2-1 duel at Mobile. That game was a turning point as he won his next eight decisions between May 30 and July 4. It seemed as if another outstanding season and a return to New York was in store, but disaster struck during a road trip to New Orleans.
Typhoid fever was a dreaded disease a century ago. George Grossart, Cozy Dolan and minor-league player-manager Julius "Hub" Knoll had all died of typhoid in the early years of the twentieth century.
However, there was no mention of typhoid or indeed any serious concern when Castleton returned to Atlanta on July 10. Within a week, the Constitution reported that he was "a very sick man...and is out of the game for the remainder of the season." He spent his twenty-third birthday in the hospital, receiving flowers from members of the Atlanta and Mobile teams.
He lost 35 pounds by the time he left the hospital after forty days. He also lost a chance to return to New York. As part of the agreement when Castleton was loaned to Atlanta as park rental, New York got the pick of the Atlanta roster. When the option was exercised in August, pitcher Russ Ford was chosen instead. At the time, Ford was 13-11, while Castleton had been 10-5 at the time of his illness.
Despite missing about half of the season, Roy was chosen on at least one Southern Association All-Star team for 1908. The team, chosen by Constitution sports editor Dick Jemison, also included past and future major leaguers Ted Breitenstein, Bill Bernhard and Tris Speaker.
Off-season reports indicated Castleton was regaining his health, and expectations were high for the 1909 season. A development in the major leagues would also affect his career. Griffith left New York to become manager of the Reds, and decided to move Cincinnati's spring training base to Atlanta.
During spring training, it seemed Roy would return to form after his serious illness. After two impressive spring starts against South Atlantic League teams, Griffith purchased an option for Castleton's services. The option required Cincinnati to acquire his contract from Atlanta on or before August 20.
But Roy struggled after the impressive start, walking five in a 4-3 loss to Birmingham on opening day.
Castleton's next start was even worse. He lost 11-7 at Nashville and also argued with catcher Sid Smith before leaving the game. Roy dropped to 0-3 before finally winning the first game of a doubleheader against Birmingham on May 1. He won that afternoon 3-0 in ten innings, allowing five hits while striking out ten and walking two.
Despite the strong outing, Castleton appeared in just one more game for Atlanta. Just days later, Roy was sent home with an illness first described as malaria and later as food poisoning. Rumors of an impending release were also circulating.
Those rumors became reality, at least in part, on May 16. Cincinnati exercised their option early and purchased the southpaw. Roy offered an insight as to why he wanted out of Atlanta: "I do not feel I can stay in shape in this climate. Atlanta is all right, and so is Nashville, but when I get to the other towns, I get sick and don't feel like working. I feel with a little rest up and a cooler climate I will be myself and able to do good work, and I hope to give Cincinnati the best there is in me." Atlanta manager Billy Smith said the purchase price was $1,500 and agreed the move was made at Castleton's request.
Upon reporting to the Reds, Castleton didn't pitch for nearly a month. He finally got into a game for the Reds on June 9. The Reds and Boston had a doubleheader scheduled, but due to rain and darkness just one game was played. An estimated 2,500 fans witnessed a 13-2 Reds win that afternoon. The Boston Globe described him as "a rather slight southpaw." The Globe was impressed with his National League debut. "True, he allowed 11 hits, but these were so nicely distributed that had not [first baseman Dick] Hoblitzel made a mess of a rap in the seventh inning, not a run would have trickled over the plate for the Doves. Castleton began the work of destructor in the first inning by fanning [Johnny] Bates and [Fred] Stem and causing [Bill] Sweeney to roll out. In the second he got [Ginger] Beaumont and [Claude] Ritchey on flies and downed [Bill] Dahlen by the three-whiff route." Castleton finished with four strikeouts and just one walk.
Unfortunately, Castleton's physical problems continued. Ring Lardner reported Castleton was in Waukesha, Wisconsin, trying to regain his health while the Reds were playing in Chicago. He soon returned to the Reds but was used infrequently. He pitched in a mop-up role in games at Pittsburgh and Boston, before entering a July 25 game versus St. Louis in the eleventh inning. The appearance was a disaster. He allowed three hits and walked three in two innings, giving up three runs in two innings of work. The loss evened his record at 1-1. As it turned out, this was his last appearance of 1909.
Brooklyn pitcher Nap Rucker had a brief vacation from his team in early August and discussed Castleton's health with a reporter from the Atlanta Constitution. "Roy does not look well, and it is doubtful if the boy ever pitches again. He told me himself that he was leaving for his home in Salt Lake City for the remainder of the season."
Castleton reported early to the Reds 1910 spring training site, Hot Springs, Arkansas. He started the season with Cincinnati, but pitched just once in the season's first month.
Ring Lardner described his inauspicious debut versus the Cubs: "Castleton went in to pitch for the Reds in the fifth. He lasted about a minute. [Ginger] Beaumont and [Frank "Wildfire"] Schulte singled and [Frank] Chance sacrificed. Steiny [Harry Steinfeldt] walked. With the bases full, Castleton uncorked a wild pitch and Beaumont loafed home. Griffith asked O'Day to call time while he switched pitchers, and Hank sent him off the field for too much talk." The two runners left on base scored, and Castleton was the losing pitcher that afternoon.
It was almost exactly a month before Griffith gave him another chance. On May 15, he made the most of a rare start, beating Nap Rucker 2-1 in a pitcher's duel. Castleton allowed five hits, walked five and struck out three. Given another start against the Giants four days later, he was knocked out in the third after surrendering four consecutive singles. Castleton would make just one more major league appearance. He entered a May 29 game against St. Louis in relief of Jack Rowan in the fifth and was the losing pitcher, working just 1 2/3 innings.
When Cincinnati started on an eastern road trip after that series, Castleton didn't go with them. Instead he went west. The Reds sold him to Los Angeles of the Pacific Coast League. Whether it was after effects of his 1908 illness or inactivity, Castleton's contribution to the 1910 Reds was minimal.
Having expressed an interest in pitching closer to home, Castleton now had the chance. In his debut at San Francisco on June 18, he struck out ten and allowed just four hits in eight innings. The Los Angeles Times said of his debut: "It is some time since a new pitcher has shown greater promise than Castleton. He has the speed of any twirler in the league, and his ball breaks in a most bewildering jumping fashion. Nearly every Seal succumbed to his curves during the afternoon, and his cool confident manner made him many friends among the crowd." Nevertheless, he lost the game 4-2.
On July 16, Castleton was even better, throwing a one-hitter against crosstown rival Vernon. He gave up a single to the second batter he faced and walked just one batter that afternoon. According to the Times, "Castleton's heaving has not been exceeded on the Coast this year. The Hooligans [Times nickname for Vernon team] knocked the ball right into the fielders' hands and everything was handled in major league fashion. Several pretty catches and stops and throws gave the fans plenty of opportunity to howl." The game was a typical deadball era pitcher's duel. Castleton won 1-0.
Roy seemed to have regained his pre-illness form, but the rest of the season was a struggle. Some of it was wildness. In a couple of late July losses, he allowed multiple walks at inopportune times leading directly to the loss. Through late August his record was just 4-7 with 65 hits and 33 walks surrendered in 66 innings.
Still, it was lack of fielding support that attracted the most comment by Castleton and by the Times. The Times account of a 3-0 September 12 loss to San Francisco summed it up: "Castleton said some time ago that every time he pitched, a lot of the boys hauled their gum boots out of the clubhouse and put them on so they could kick the ball around. Whether or not this is true is another story, but the fact remains that the Angels threw the pill around in a weird way at times yesterday and the Seals got two runs off the weirdness." The Times commented on poor defensive support for Roy several other times that season, even offering the opinion that the poor defense discouraged him and diminished his effectiveness. Whatever the cause, he finished 1910 with a less than impressive 8-15 record with 95 strikeouts and 68 walks.
Though originally expected to remain with the Angels for the 1911 season, the Los Angeles manager asked for waivers on Castleton. Hap Hogan, manager of crosstown rival Vernon, claimed him in early February.
Castleton opened the 1911 season with Vernon as one of seven pitchers on the staff. After an undistinguished relief appearance, he started his first game of the season at San Francisco on April 1. He shut the Seals out 4-0 that afternoon and received positive comment in the Los Angeles Times: "San Francisco players were bowled over with the rapidity of a crack rolling a big ball down the alley. A couple of their men did reach second, but Castleton scattered his five hits over as many innings, and there was no question but that the Seals were hopelessly distanced. Castleton, for a southpaw, was remarkably steady, allowing but one walk."
He won his next start 4-2 at home against Portland on April 6, shaky control working to his advantage. The Times described him as "fearfully and wonderfully wild at times, and, in addition to three bases on balls, he made two wild pitches. The fact that he didn't bounce the ball off any solid-ivory heads shows that the Beavers were pretty much there in the ducking line."
After a couple of losses, Castleton finished April 1911 with a flourish. He beat Oakland twice in the same series, one by a shutout, but he saved his best effort for his last start of the month. On April 28, Vernon faced the crosstown rival Angels, and Roy pitched his third shutout of the month. The Times said of the game, "what he did to them [the Angels] was what John D. Rockefeller does to all of us. He simply hung out the sign that it was his busy day and consequently there were no transactions in the run line. Castleton just made monkeys out of them. [He] curved them around their necks, fanned seven of them and those who could hit the pill out of the infield thought they were lucky." He allowed just five hits while improving his record to 5-2.
Roy struggled during May and June. He lasted two and one innings, respectively, in a pair of starts at Sacramento. He also was knocked out after two innings in a mid June start against Portland. The Times offered colorful comment on the game: "Castleton thought early yesterday morning that he understands something about heaving, but early yesterday afternoon he didn't think so. In fact, he did not think much of himself. If he did think it couldn't have been much, for the champions slammed him for eight runs and six hits in the first two innings. This was naturally all he wanted, and he was glad to get under the bench where they could not even see him." His control was abysmal that day. He walked five batters in the first inning.
The team's performance was as inconsistent as Castleton's. The fortunes of the team and the pitcher began to change during a July road trip. He shut out San Francisco and Portland in consecutive starts and won another game in the series at Portland with his arm and his bat. He allowed just four hits in ten innings and singled in the winning run. He won seven consecutive starts in July and August, allowing five or fewer hits in four of the games. During the streak, Vernon took the lead in a tight pennant race.
By mid-September the Tigers were in Portland facing their top competition for the PCL pennant. His manager hoped to use him twice in the key series, but it rained on Tuesday, September 14. Castleton slipped on wet pavement that afternoon and sprained an ankle. He pitched in the first game of a doubleheader three days later, but seemed to be affected by the injury. Poor defense and six walks helped the Beavers to a 5-4 win, and a win in the second game gave them a decided advantage with a month left in the season.
Castleton recovered quickly from the injury, shutting out Sacramento and allowing just one run on five hits against Los Angeles. The latter win tied Vernon for the league lead, but fatigue was evident. He didn't win another game, and was hit hard in two October starts. He finished the season with a 22-13 record, striking out 146 and walking 76 in 327 innings. The heavy workload would take its toll.
Though there was a brief misunderstanding about his contract, Castleton soon signed with Vernon for the 1912 season. According to Roy, "There was some mistake in my address and Hap's letter containing the contract lay in the post office unclaimed since January 5. As I did not hear from Hogan by February 1 I thought I was free, according to the new rules. Since February 5 I have heard from Hogan and believe everything is all right between us." He spent the off-season working in the auditor's department of the Oregon Short Line Railroad and kept in shape by shoveling snow and playing hockey.
After reporting for spring training, Roy began the season well. He pitched impressively in exhibition games, including one against the University of Southern California, and was counted on to play a major role in Vernon's hopes for the pennant.
Castleton struggled in April, but seemed to return to form in May. After a victory over Oakland on May 12, the Times said, "A double hop that Castleton put on what he had yesterday sewed the Oaks up into knots and easily won him a ball game with apparently little effort on his part." Castleton won six out of seven starts, completing them all, between April 28 and May 26.
It seemed he was on the verge of another outstanding season, but the innings were taking their toll. Castleton missed a month with what was described as a wrenched back and injured pitching arm.
Returning in late June, Castleton won three straight decisions but was evidently still pitching in pain. Of one of those wins, the Times said, "Portland secured enough hits [ten] off Castleton to win an ordinary game, but this much must be said for the southpaw, he breezed along under half speed until forced to show his colors then tightened up short."
By the beginning of August even this didn't work. Roy was seldom able to finish what he started. Since Vernon was in another pennant race, he continued to take the mound, and resented being removed. After he was taken out in the tenth inning of an eleven-inning loss to Los Angeles, the Times described Castleton as "about as sore as a half-boiled owl when he was yanked, for he had pitched a fine game up to that time."
Pitching in pain, Roy soon took his frustration out on Los Angeles Times reporter Grey Oliver. After being removed in the fifth inning against Oakland On August 16, he said, "What do you think. Me going out there with a sore shoulder and then them stickin' Stinson in right field with a broken leg and not able to catch fly balls that any man would have caught if he was right. Guess if some of those fly balls had been caught there would not have been so many hits. And then you fellows say we are knocked out of the box."
Castleton completed just one game after August 1, 1912, but it was a masterpiece. On September 4, he won a key game against Vernon's top competitor, Oakland. He shut the Oaks out on four hits. The Times said: "Castleton pitched ball like a champion. He had all those Oaks going south and they knew it. He never was in danger. Up to the eighth round, only two drives had been registered against the left hander. Two more were bunched in the eighth, but this did not do any good, for as soon as they were made, Castleton tightened up and the men behind him made short work of the struggling Oaklanders. They were simply outclassed."
Henry Heitmuller, the PCL's leading hitter, died of typhoid fever on October 8. The death must have affected Roy as much as it did Heitmuller's teammates. It almost certainly brought back memories of his own struggle with the illness and probably had a major impact on a decision he'd have to make a few months later. The day of the funeral in San Francisco, the Los Angeles- Vernon game was stopped for ten minutes in the third inning. Castleton pitched and won that day though he was removed in the eighth inning. That game was his final win of the season--and his career. He had one final appearance after that, saving a win in the first game of the Tigers season ending doubleheader win over Portland.
He finished 1912 with a 13-8 record in 222 innings striking out 92 and walking 83. He hit .152 in 79 at-bats. After the season he reportedly toured Australia with a PCL all-star team organized by J. Cal Ewing of Oakland.
Shortly after the season it became clear that Castleton wouldn't return to Vernon. A rumored trade to Portland fell through, and in late January of 1913, he was sold to Nashville of the Southern Association. An early report said he'd signed with the team, but later news indicated he was holding out. An April 7 article in the Atlanta Constitution said, "while in the shadow of Brigham Young's temple, out in Utah, Roy Castleton is wailing to be shipped to some climate where the typhus germ is an alien." Perhaps if the Nashville ballpark hadn't been flooded and Vols manager Bill Schwartz hadn't been stranded in Ohio, also due to flood, a deal might have been worked out. It wasn't and Roy Castleton retired from baseball.
Despite bad memories of Southern baseball, Roy was still mentioned among the best to play in Atlanta and the Southern association into the 1920s on the strength of his success in 1907 and the first half of 1908.
On July 9, 1918, Roy married 25-year-old Esther Kelson in Salt Lake City. At the time he was working as an accountant for Scott & Hadley, a stock brokerage firm in Salt Lake City. Roy, Esther and several of his siblings continued to live with his parents at least until Charles Castleton Sr. died in 1922.
Roy and Esther had no children and later moved to Los Angeles, where he worked as a water heater inspector. Suffering from diabetes, Roy died there of a cerebral hemorrhage on June 24, 1967, and was returned to Salt Lake City for burial in the Salt Lake City Cemetery.
Ogden (Utah) Standard Examiner, 1904-08.
Youngstown (Ohio) Telegram, 1905-06.
Akron (Ohio) Beacon Journal, 1905,06,09.
Lancaster (Ohio) Gazette, 1906
Lancaster (Ohio) Eagle, 1906.
Newark (Ohio) American Tribune, 1906.
Dayton (Ohio) Journal, 1906,09,10.
Toledo (Ohio) Blade, 1902.
New York Times, 1907, 1910.
New York Mail, 1907.
New York Globe, 1907.
Atlanta Journal, 1907.
Atlanta Constitution, 1907-09, 1912-13, 1921.
Chicago Tribune, 1907, 1909-10.
Boston Globe, 1909-10.
Washington Post, 1910.
Los Angeles Times, 1910-12.
Salt Lake City Tribune, 1967, 1987.
Other Baseball Sources
SABR Online Encyclopedia.
Wright, Marshall, Southern Association in Baseball, 1885-1961 (McFarland).
Lee, Bill, Baseball Necrology (McFarland).
Jones, Kevin, e-mails December 2005.
Nelson, Rod, e-mail July 2005.
Wendt, Paul (forwarded by Kevin Jones) e-mail December 2005.
Utah, 1870, 1880, 1900, 1910, 1920.
Other Genealogical Sources
Passenger Lists Port of New York.
World War I Draft Registration Card Royal Eugene Castleton.
Utah Pioneers and Prominent Men.
Ancestor Histories Salt Lake City Chapter Sons Of Utah Pioneers.
Other Online Sources
Jensen, Richard L., "Immigration to Utah," (Utah Historical Encyclopedia).
Kimball, Stanley, "Mormon Trail in Utah," (Utah Historical Encyclopedia).
McCormick, John S.," Salt Lake City" (Utah Historical Encyclopedia). |
Buddhism and Christianity
by John B. Cobb, Jr.
John B. Cobb, Jr., Ph.D. is Professor of Theology Emeritus at the Claremont School of Theology, Claremont, California, and Co-Director of the Center for Process Studies there. His many books currently in print include: Reclaiming the Church (1997); with Herman Daly, For the Common Good; Becoming a Thinking Christian (1993); Sustainability (1992); Can Christ Become Good News Again? (1991); ed. with Christopher Ives, The Emptying God: a Buddhist-Jewish-Christian Conversation (1990); with Charles Birch, The Liberation of Life; and with David Griffin, Process Theology: An Introductory Exposition (1977). He is a retired minister in the United Methodist Church. His email address is email@example.com.. This essay was one of two lectures given at Bangor Theological Seminary, January 26-27,2004. The second essay is entitled "Beyond Pluralism." This material was prepared for Religion Online by Ted and Winnie Brock.
In my previous lecture I talked about the need to consider separately every other religious tradition and how as Christians we should understand and relate to each. We must consider the history of our past relation to it, its strengths and weaknesses, and the practical effects of taking particular actions. We may end up with some generalizations, but we should move from the particular to the general, not from ideas about religion in general to the understanding of and response to particular traditions.
In that lecture I talked first about what we could contribute to others, and then about what we might learn from them. But I also noted that, at least in today's world, after centuries in which we talked far more than we listened, it is time to put listening first. After we have learned from the wisdom of the other tradition and been transformed by what we learn, we are in much better position to be heard when we speak.
Although I illustrated my point of the need to accept diversity in many ways with reference to Judaism, Hinduism, and Shinto, and spoke of what we might offer them and learn from them, all my comments were very brief. The reality, on the other hand, is that we need to wrestle with these questions profoundly and extensively, much as the Church Fathers wrestled with the wisdom of the Greek tradition. This is a major task for theology in this twenty-first century. It was begun in the twentieth, but only begun. It will require the joint work of many people.
My own small contribution to initiating this work has been in relation to Buddhism. There are several reasons that I have focused on the relation of Christianity and Buddhism.
First, because I lived in Japan as a boy, I have had more contacts there than in countries where other religious traditions prevail. Religiously thoughtful Japanese are typically interested in Buddhism, and many of them are devoted Buddhists. I have had opportunities to talk with a number of them. From there our Buddhist-Christian dialogues spread to include Buddhists from other regions, especially Tibet and Thailand.
Second, among my mentors, teachers, and friends several were especially interested in Buddhism. This was true of both Whitehead and Charles Hartshorne, and also of Thomas Altizer. This interest reflected not only the specifics of their religious concerns, but also a wider cultural phenomenon. Buddhism has fascinated Christians for centuries.
Third, one reason that Whitehead and Hartshorne were interested in Buddhism also affects me directly. In a remarkable way, such Buddhist thinkers as Nagarjuna anticipated the insights of the contemporary process philosophical movement. Western process thought did not derive from Buddhism, but it cannot but recognize that it adopted its views for some of the same reasons two millennia later.
Fourth, whereas Westerners have come to these philosophical conclusions rather recently, Buddhists have lived with them for many centuries. Whereas Westerners have done so in a context in which philosophy and religion are considered quite distinct, Buddhists have lived with them in a context where no such walls of separation existed. For them, the questions are: What are the existential or spiritual implications of these insights? How can they support the quest for enlightenment? Apart from the influence of Buddhism, Western process thinkers have hardly asked these questions, much less answered them. They have obvious importance for a theologian.
It is evident that my judgments as to what Christians can learn from Buddhists are greatly affected by the congeniality of its basic insight with my own beliefs. That means that the judgments I will offer in this lecture are more directly dependent on my philosophic views than was the case in my previous lecture. Yet I would not like for those of you who do not want to be dependent on a particular philosophy to shut me off. Accordingly, I will discuss the problem in Christian theology to which process theology gives an answer first, before describing the answer given in the Whiteheadian process tradition. Many theologians not committed to the specificities of the answer of process theology agree that Christianity has had a problem with substance thought.
In my previous lecture I expressed my admiration for the work of the Church Fathers in the Hellenization of Christianity. Without this indigenization of Christianity in the Greco-Roman world, the movement would have failed. But, even so, Christianity has paid a high price. The Fathers did not abandon the biblical story in order to make the faith understandable to Greeks, but they did impose a conceptuality in deep tension with it.
The Bible is chiefly written in story form. That is, it is about events. It tells what people did. One can tease out of the way the story is told some ideas about the structure of human beings: body, emotions, will, soul, spirit, and so forth. But no biblical author attempts to describe this anthropology. We can say more, on the basis of the Bible, about relationships among people than about the nature of the people who are related.
Similarly the Bible speaks a great deal about what God does, and it provides many images and descriptive adjectives. But it tells us almost nothing about the nature of the divine existence in itself. Since the stories are told over a period of a thousand years or more, their depictions of the divine character vary. It is very difficult to reconcile some of the stories with some of what is said about God's character.
Sophisticated Greek audiences felt these gaps keenly. They could not believe that God acted in the barbaric way portrayed in some stories. Other stories presented God as all too human. In general the questions the Greeks asked required answers at what we might call the ontological level. What is the nature of God? How does God affect what happens in the world? Their answers required them also to judge that many of the stories should not be taken at face value. They decided that they were in the text as a source of moral and theological lessons rather than in order to satisfy curiosity about ancient events. Broadly speaking cosmology replaced historical narrative as primary. The anthropomorphic God of the Bible was replaced by a new doctrine that wove together elements from the Bible and from Greek philosophy in a new synthesis.
The new synthesis replaced temporality and history with a nontemporal eternality. For the Greeks, to be divine was to be eternal in the specific sense of being above or beyond time. One of the major attributes of God became immutability. In the Bible we are told that God is faithful to God's promises and that God's character never changes. But this new synthesis went far beyond that. Strictly speaking what happens in the world cannot make any difference to a God who cannot change in any way.
The only alternative to this denial that God know what happens in the world was to say that time is unreal for God, so that the whole course of events affects God eternally. Even when time is sacrificed in order to acknowledge God's knowledge of the world, God was not really allowed to care. That would introduce negative feelings into God, and that was thought to be incompatible with God's perfection. This whole pattern of thinking was antithetical to scripture. For the Bible, we are called to serve and please God. But the new doctrine of God renders that idea meaningless.
Another problem came from the new doctrine of divine omnipotence. Many suppose that this is biblical. It is not. Certainly the Bible speaks a great deal about God's amazing power, but it does not deny some lesser degree of power to other spiritual entities and to human beings and other animals. God is the most powerful being, but there is no reason to say that God has all the power. One problem was that, in the Septuagint, "Cosmocrator", ruler of the cosmos, replaced Shaddai, the proper name for God in part of the Pentateuch and in Job. In translation into Western languages, beginning with Latin, this was again mistranslated as almighty. Cosmocrator means ruler over all, but it does not deny power to those who are ruled. The Cosmocrator has to be very powerful to control the powerful forces within the cosmos. To say the ruler of the cosmos rules over powerless beings actually reduces, even dissolves, the affirmation of God's power. Also, Christianity has suffered immensely from having to explain how there can be evil in a world in which a good God has all the power. Formulated in that quite non-biblical way, the question has no possible answer, as the long history of theodicy makes quite clear. The Bible tells a dramatic story of many actors in which God plays the primary and ultimate role. If God is the only actor, there is no real drama.
Of the many other problems introduced into Christian thought by the relation of the new synthesis to scripture, I will mention only one more. In the Bible, boundaries are somewhat fluid. This is important for understanding the way God acts in the world and the human experience of God's presence. Paul is particularly interesting here. He speaks of our participation in Jesus' faithfulness, suffering, death, and burial. He speaks of Christ in us and our being in Christ. He speaks of the indwelling of the Spirit. He speaks of how we in the church are members one of another.
Substance thought cannot make sense of any of this. Instead of the Spirit working righteousness in the hearts of the faithful, the interpreters could speak only of God's acquitting believers of sin after the manner of a judge. Instead of the faithful participating in the life and death of Jesus, they can at best imitate. Instead of the faithful being 'en Christo" at best, God treats Jesus as a substitute offering.
Equally important for the development of Christian theology was the idea of the incarnation derived from the prologue of the Gospel of John. Here the Word of God is said to become flesh. That this had occurred became the hallmark of orthodox Christology. Just what it means became the topic of generations of debate and mutual recriminations. I am an enthusiast for the idea of God's incarnation in the world and especially in Jesus, but you will not be surprised to hear me say that the idea became an acute problem because of the Hellenistic categories in which it was discussed.
If the discussion had continued in biblical terms, ontological precision would not have been sought. The Bible talks often enough about the Word coming to someone, and sometimes of its operation within them. Similarly God's presence in the world can be spoken of in terms of God's Spirit, God's Wisdom, and God's glory. That these are present from time to time in people is part of biblical understanding. That such presence is affirmed of Jesus is to be expected. That it is affirmed in stronger language of Jesus than of others does not surprise or occasion intellectual puzzlement. The language of Antioch, where biblical images played a larger role than in Alexandria, was of the Word indwelling Jesus.
But for those who could think only in Greek philosophical categories, this was not clear or sufficient. For them the world was made up of entities that were identified as ousia. This was translated into Latin as substantia which, of course, becomes substance in English. There were philosophical debates among Greek philosophers as to exactly what this word means, but common to their usage was the idea expressed in the adage, "two substances cannot occupy the same space at the same time." If one takes this for granted as a fundamental principle, one entity, even a divine one, cannot indwell another. The incarnation becomes incomprehensible. The doctrine requires that Jesus be both fully human and fully God, but God and humanity cannot occupy the same space at the same time.
Those who were strongly shaped by this conceptuality proposed that some feature of Jesus' humanity was replaced by the presence of the Word. Those who were more concerned with the biblical story, and those whose theology required Jesus' full humanity, insisted that Jesus was not lacking in any human feature. The creeds ended up in paradox and stalemate. But later, after the time of the creeds was passed, it came to be considered orthodox to think that Jesus selfhood or person or "I" was only divine. This is surely a wholly unbiblical idea developed because of substance thinking.
These were by no means the only places where Greek philosophy, based on the ontology of ousia, blocked the more natural expression of biblical ideas, but they should suffice to indicate the problem. Even though the Reformers tried to liberate the Bible from philosophical thinking, they hardly began to deal with the deep hold of substance thought on theology. Sadly, if one is unwilling to think philosophically, one will be the servant of its existing form.
Modern theology was founded with an even sharper focus on substances than had characterized medieval thought. Perhaps just because it was so explicit and so emphasized, it came under serious question for the first time in Western history. It gradually became clear that when substances were clearly distinguished from their changing attributes they could not be thought at all. Finally, Hume was brave enough to reject them out of hand. Kant tried to restore substance as a necessary category for human thought, but for the most part Western philosophy has simply bypassed the discussion. It learned from Kant that we should abandon metaphysics and think instead about human thinking.
This opens the door for the reappropriation of more biblical modes of thought. For example, there has been much talk of narrative theology. Others have reduced theology to language. If language is no longer thought to refer to something beside itself, we are free to use biblical language quite uncritically. However, we should recognize that when we do so, we are speaking quite differently than the biblical authors who naively thought they were talking about actual events.
I have summarized the story quite independently of the process conclusion that I myself draw from it. I hope that it makes clear why the encounter of Christianity with Buddhism is important. It is a movement of thought and spiritual life that for thousands of years has rejected the idea of substance. For the West, thus far, the recognition that there are no substances has led primarily to the idea that talking about what is or what is not is a mistake. If there are no substances, it is assumed, then there is no reality beyond our experience or our language. There is no meaning in asking about the existence or nature of God. We can only talk about our symbols, our language, or our experience of God. This abandonment of realism by the educated elite has driven a deep wedge between a great deal of theology, on the one side, and the actual piety of the church on the other. This is not a healthy situation.
There have been many philosophical and theological responses. One family of these responses undertakes to continue the discussion of reality by replacing substances with events. There are no substances, but this does not mean that nothing happens. It simply means that the events are the reality. There is not another type of entity underlying the events and acting through them. I subscribe to this view. After trying so long to think of events as simply the product of matter in motion, it is time to think of matter as a pattern of events. In a very general sense this is the triumph of biblical historical thinking over Greek ontological thinking. But it is not really that, because it asks and answers many of the same questions that constituted Greek philosophy. Instead of abandoning ontology or metaphysics, it proposes a different one.
Sadly, from my point of view, this option is viewed with great suspicion in intellectual circles. Hume and Kant continue to shape the discussion. They emphasize what cannot be done, and what cannot be done is to develop a new metaphysics. They do not find it necessary to criticize the new philosophy in detail, since it is the enterprise itself they reject. It is, from the dominant point of view naively realistic and unrealistically ambitious. The time for such grand schemes has, in their perspective, long past. They can even appeal to Buddhism in support of this rejection of the constructive enterprise.
Nevertheless, we do encounter in Buddhism an alternative to substance thinking that does not simply bypass the issue of what is. And this introduces Buddhism, or more precisely multiple Buddhist voices, into the present discussion as fruitful contributors. I will make no effort to describe the diversity within Buddhism, which is vast. I focus only one doctrine that is widely held. In place of substances, Buddhists speak of pratitya samutpada. This is translated as dependent origination. This means that what there is is always a coming into be out of what is other. Therefore, nothing exists in itself, as a substance is thought to do. Each thing is what it is in any moment in derivation from what it is not.
Buddhists illustrate this in many ways. I will illustrate it in the way that I find most convincing. I can say, at least, that when I have used this illustration, Buddhists have not objected. Substance thinking has typically asked us to think about a stone, or a tree, or a star. When we begin with these sorts of entities and generalize about them we are likely to end up with a metaphysics of substance, or, if we dissolve the substance as Hume did so brilliantly, we end up with nothing but our own sense data. Ultimately we must be able to show that these sorts of entities are also instances of dependent origination, but when Buddhists begin with them, it seems to me, they are often not so convincing.
My example, then, is a moment of human experience. What is that? It is quite apparently something than happens. In itself one is not tempted to think of it as a substance. If one insists on substance thinking, one will posit a subject who enjoys such experience. Ordinary language encourages this. One says, I see the dog; so there is a subject separate from the experience who might be seeing a cat instead. My point is that one can begin with the event and then be led by ordinary English language to posit a substantial subject and a substantial object outside the experiential event.
However, once we have recognized that analysis of the given in terms of postulated substances does not work, we may be willing to simply examine the experience itself. Of what does it consist? I will identify only a few elements. There are feelings of bodily events, perhaps an aching back. There are memories of past events, perhaps a delicious meal recently enjoyed. There are feelings about other people, perhaps a beloved child who is in danger. There are relations to the environment that express themselves in colors and sounds.
Now comes a further question. Is there first of all an experience that then relates to the body, the past, other people and, the physical environment in these ways? That would introduce a new substance, the experience as such. But that analysis fails. We find in experience no experience as such. The only experience there is is the experience of other things. The experience in its concrete actuality is the togetherness of these other things. The aching back, the tasty meal, and so forth originate the experience. It is an instance of dependent origination. This momentary experience ceases to be and shares in giving rise to a successor experience.
Attending to human experience is important for Buddhists, because central to their teaching is the idea of no-self. This means, what we have already noted, that there is no substance underlying the flow of experiential events. The events themselves are all that there is. This does not mean that there are not other events, in the brain, for example, that contribute to the dependent origination of human experience. But there is no self to be distinguished from the events.
Many Christians find this disturbing and offensive. They feel that human beings are reduced or demeaned by this idea. Buddhists do not think so. They find the idea in principle liberating, and they hope to realize, existentially, its truth. If they doe so, they are convinced, they can attain a serenity that it otherwise eludes human beings.
The Buddhist doctrine of no-self is no doubt in tension with Christian thought, but it was not developed to counter it. It must be understood in its Indian context. We noted in the previous lecture that a major strand of Hindu thought celebrates the identity of Atman with Brahman. Here Atman is the human self and Brahman is ultimate reality. In the Hindu context, these are typically, although not always, thought of as substances.
This standard Hindu analysis moves in the opposite direction from the Buddhist. When the Hindu reflects about experience, it dissolves into the world of appearances. This is much the same as happened to Hume. But Vedantists are convinced that the reality of the Atman is not affected by the play of appearances. The real self is not the empirical self that appears in that play. The reality underlying human experience is timeless and without differentiating qualities. Similarly the reality that underlies the play of appearances objectively, which means the underlying reality of all things, is also without any differentiating qualities. It is Brahman, the ultimate. It could also be called Being Itself. This Being Itself is the Being Itself of all that is including the self. Hence Atman and Brahman are one.
Agreeing to this analysis may have some intrinsic value, but for the Vedantist the task is to realize its truth existentially or mystically. Particular yogic disciplines are designed to do so. There is no doubt that they can lead to extraordinary states of consciousness.
Nevertheless, it is this analysis against which Buddhists have argued. They do so, as we have seen, in a purely theoretical way, denying that Atman or Brahman underlies the flux of events. They also see this theory as leading people to depreciate that actual events of the world viewing them as appearance and the appearance as illusory. The Buddhist seeks to realize the existential meaning of the fact that there is no Atman rather than of the unity of Atman and Brahman.
Many Hindus think that Buddhist exaggerate their difference from this form of Hinduism. It is possible to think of Atman and Brahman as dynamic such that they name an always coming into being rather than a static underlying Being. Then the realization of the identity of Atman and Brahman is not so different from the recognition that every event including the event of human experience is an instance of dependent origination. In some such way it may be possible to get past the sheer contradiction that seems, from the Buddhist side, to be entailed.
Nevertheless, there are important differences between typical Buddhist and Vedantic practices. I like to tell about an experiment I read of many years ago. Sadly I do not now have the report, so that my description is based on failing memory. Nevertheless, I believe the basic point is safe.
At a Catholic university in Tokyo, experimenters hooked up their subjects to instrument measuring brain activity. There were three groups involved. The first were people who prayed or meditated with no specially developed discipline. The second was a group of yogic practitioners. The third group was composed of practitioners of Zen disciplines.
All were asked to meditate or pray. After a few minutes the experimenter sounded a raucous buzzer. Then after a pause repeated this several times. Afterwards the brain waves of the different participants were studied.
The first group responded in the expected way. The unpleasant sound interrupting their prayers led to significantly heightened activity the first time. Its effect was less the second time and continued to diminish as it was repeated.
The yogi practitioners, on the other hand did not respond at all. Presumably they had shut out the external world and moved deeply into their interior. What happened externally did not matter to them.
The Buddhists, however, responded moderately to the first buzzer and similarly to each succeeding one. I want to use this response as a way of clarifying the implications of the no-self doctrine.
Buddhists judge that as long we understand our experience as being grounded in enduring selves, we will live in our past and future. We have regrets or feel pride in past acts and anxieties about what the future will bring. When we break with this erroneous view, we can live moment by moment in the present. Living in this way, we can be fully present to just what is, not interpreting it in terms of its harmfulness or benefit to ourselves, but letting it be. We experience it and then let it go. If it happens again, we will again experience it, and then let it go.
Buddhist meditational discipline, therefore, does not separate practitioners from ordinary life. It enables them to live that life in full immediacy moment by moment being fully present to whatever or whoever is there. There is no denial or pain or suffering, presumably the buzzer was unpleasant. But the response to the unpleasant sound was simply the recognition of the fact that it occurred. It was not viewed as an interruption of something else that was more important or a negative portent of what would happen in the future,
There is surely much in this that is interesting and attractive to Christians. If letting go of self leads to this kind of serenity, one need not dread it. It is not clear what positive role this self would play in Christianity.
I am not trying to say here just how Christians should respond to the few features of Buddhism I have highlighted. I hope I have said enough to show that one cannot quickly say that we have nothing to learn or that appropriating anything from Buddhism would be contrary to our faithfulness to Jesus Christ. We may on reflection decide that we cannot agree with the ideas or consider the goal the right one for Christians. But even if this is the decision, Christians who have come to it in response to Buddhism will be changed in the process. We cannot deal seriously with new issues, whatever the outcome, without being broadened. If we reaffirm our inherited answers, we now understand differently the questions to which they respond.
My own judgment, as I have said, is that Buddhists are basically correct in their understanding of reality, and that we can learn from them the positive consequences of realizing this. This is a major contribution that Buddhists can make to us. It will take us a long time truly to learn it and be reshaped by it. Until we have learned this, it will be difficult for them to listen to us.
To explain this let me talk about what I think we have to teach. Most Buddhists have no place for God in their thought. Originally, they denied the reality of Brahman, as of Atman. They did not deny that their were superhuman beings who were worshipped by some people, but they taught that that was a mistake. What was truly important was enlightenment and these gods were distractions. At that time they knew nothing of the God of the Abrahamic traditions; so their denial was not directed at this kind of monotheism. However, when they encountered the belief in this deity, either in its Muslim or its Christian form, they rejected it too. Let us consider why.
My previous discussion should provide at least part of the answer. An all powerful deity is completely incompatible with an understanding of dependent origination. This doctrine requires that everything be partly determined by everything else. That attributes power to all things. Further the idea of substance played a large role in the doctrine of God. Indeed, God was generally viewed as the ideal embodiment of substance, totally self-contained and self-sufficient, incapable of being affected by anything. Such a way of thinking closely associated the God of Abraham with Brahman.
The only kind of deity a Buddhist could seriously consider would be one who is an instance of pratitya samutpada. Such a deity would be arising dependently from all other entities and would be part of that from which all other entities arise dependently. Although many Buddhists have no interest in speculating about such a deity, their basic vision of reality is not undercut by it.
Now we must ask whether such a deity could be understood as the one whom Jesus called Abba. My view is that it would fit much better than the traditional divine substance. In Jesus' understanding there is interaction between God and the world. God cares what happens in the world. Indeed there are indications that he thought that God knew every event in the world. Since the biblical understanding of knowing goes far beyond the merely cognitive we might say that every event in the world affects God, even participates in constituting the divine experience. He certainly thought that God acted in the world, and it would not be upsetting of his thought for God to be said to be participating in or informing every event whatsoever. A Buddhized formulation would be close to the one Jesus called Abba.
On the other side, despite the desire of many Buddhists to avoid any entanglement with the idea of God, there are developments in Buddhism that suggest an openness to the kind of deity of which I have spoken. Many Buddhists long for compassion and help. They turn to Buddhas for this. These Buddhas are not just enlightened human beings. They may be understood to continue to exist and to act mercifully in the world. To them can be attributed a cosmic role.
Of course the very fact that there are many Buddhas of this sort limits the similarity. But in Pure Land Buddhism, one Buddha, Amitabha is depicted aas encompassing the work of all the others. Amitabha acts cosmically for all.
True, according to the myth this is a story of one who became a Buddha after many lifetimes. It is not about an everlasting being essential to the operation of the cosmos. It is not yet God in the sense an Abrahamic monotheistic could recognize or allow.
Still, another step is possible for some Buddhists. There is a widely held Buddhist Trinity. It consists of the three bodies of Buddha. One body, the Nirmanakaya, is the manifest body, that is, if one wishes, the incarnation. It is Gotama, and, in principle anyone else who attains enlightenment. What is incarnated is the universal principle of pratitya samutpada, here called the Dharmakaya. It is beyond the distinction of good and evil. Hence it is more like Being Itself or Brahman, than like the Abrahamic God. But then there is also the Sambhogakaya, the Dharmakaya experienced by human beings as compassionate – we might say pure grace. This body of Buddha does not need to be understood as having come into being at some point. It can be seen as cosmic and everlasting. It provides a point of contact within Buddhist thought for speaking of the Christian God.
One can begin my noting the similarity. Buddhists rightly identify compassion as central to the divine. Christians know this compassion in and through Jesus Christ, but Buddhists have found it through the Buddhas. In that case, what is there left for Christians to teach.
I believe that we know God in a fuller way. Just as Buddhists can teach us a great deal about the existential meaning of realizing that there are no substances, that all things are dependently arising, so Christians can teach about the existential meaning of the fact that at the heart of the universe is compassionate understanding. Buddhists draw from this comfort and assurance. This is surely important, and Christians share in it. But we are surprised that there is so little talk of reciprocity. The question is entirely what the deity is doing and will do for us. There is virtually nothing about how we should respond beyond faithful acceptance.
For Christians God's gracious work within us is also directive. God not only comforts us, God also calls us. Buddhists clearly are often in fact responding to God's call. But there is an advantage in identifying what we are doing. Buddhists can learn intentionally to listen to the call of God.
Further, we believe that the God who comforts and calls also deserves our supreme loyalty, a loyalty that transcends that to any creaturely reality. This note of loyalty or commitment is lacking in Buddhism. In my previous lecture I spoke of how state Shinto promotes the worship of the emperor and how the Japanese Christians oppose this. They have little support from Buddhists. Buddhists as Buddhists do not promote ethnic nationalism, but they do not understand their religious thought and practice to deal with issues of this sort. In response to divine compassion they may recognize that they should place the love of the divine lover and the extension of that love to all others above national loyalty. The can come to see the worship of the divine lover as contrary to the worship of one people through that of their leader.
Buddhists are often averse to talk about justice. To them it sounds like repaying evil with evil. Alternately it may involve social engineering to make all people equal. Their interest is far more in social harmony. Nevertheless, we Christians believe that what we mean by justice has an important place in the response to God's love. There has been much that has been damaging in our quest for justice, and on this point our Buddhist partners can be very helpful. But it is quite possible that this conversation cam be a two-way street. Already, partly in response to dialogue with Christians, Sulak Sivaraksa had organized a successful movement of socially-engaged Buddhists. They are teaching Christians a great deal about how to be engaged in compassionate work without treating those with whom we disagree as enemies.
In the Jewish scriptures, the response is primarily in terms of obedience to the law given by God. Buddhists for the most part live in societies that are bound my many rules of contact. In their religious practice they seek something quite different. Christianity offers what is much more acceptable. As recipients of God's love we are called to love God and our neighbors, especially the least of these. This is not alien to Buddhism.
I have probably gone too far in suggesting just how God can come to play a part and then an enlarged part in Buddhist thought and life. The fact is that we cannot control how our message is perceived and what in it will stimulate rethinking on the part of the other. But much of what I have said comes close to my experience thus far. The steps I have described are not impossible.
The need for more careful thought about social responsibility has not expressed itself only the establishment of the movement of engaged Buddhists. I have had the ;privilege of supervising two dissertations by Buddhists on Buddhist approaches to social ethics. One was written by a Zen Buddhist. The other by a Pure Land Buddhist. They differ markedly from each other as each rooted the proposed social ethics inm the specifics of his tradition. I should say that neither followed the simple pattern I have suggested here.
One of my Pure Land Buddhists friends has become an avowed theist, much along the lines I have outlined. He sees no conflict between Buddhist teaching as understood in this tradition and the further development of its theistic tendencies. He agrees with me that Buddhism in principle is open to viewing Jesus Christ as a Buddha who revealed and taught the compassion of God.
If Christianity can give up its lingering commitment to substance thought and learn the existential meaning of the Buddhist doctrine of no-self, it will be creatively transformed. If Buddhists can recognize that the compassion they crave and experience is from God, and that the compassionate God is also one who calls us to compassionate action and to whom ultimate loyalty is due, Buddhism will be transformed. I believe that such transformation is already taking place at the periphery of each tradition. During this century it could go much further ant change the nature of the two traditions and of their relations. Taking part in this is an exciting opportunity and challenge. |
30 No. 4
Mitigating Arsenic Pollution: Bridging the Gap Between Knowledge and Practice*
by Hemda Garelick and Huw Jones
Access to uncontaminated water in sufficient quantities may be the most important requirement for healthy human societies. However, the relationship between the water supply in developing countries and the health of citizens is complex since the relationship is dependent on the provision of both appropriate quantities and quality of water. Attempts to understand this relationship began in the 1970s with the “Bradley Classification” of water-related infection. This was followed by efforts to predict the effect of variations in water quality and supply on morbidity and mortality, particularly of children under five (Cairncross and Valdmanis, 2006). Among the indicators of water quality, the feco-oral (specifically infectious diarrhea) group has historically been important and remains one of the largest water-related contributors to disease on a global scale.
According to some current estimates, approximately 1.8 million people, 90 percent of whom are children under 5, still die every year from diarrheal diseases—mostly in developing countries (WHO, 2008). It is estimated that in industrialized countries, 60 percent of diarrheal disease is attributed to unsafe water, sanitation, and hygiene, whereas in developing countries as much as 85-90 percent of diarrheal illness can be attributed to these causes (Keusch et al., 2006).
While the microbiological quality of water is directly related to human health impacts, access to appropriate quantities of water is an additional health benefit that should not be underestimated. Adequate quantities of water enable people to have improved hygiene, leading to a reduced infectious burden (Cairncross and Valdmanis, 2006). In addition, water availability beyond that required for sustenance and hygiene can confer further health benefits by allowing increased food production in the dry season.
The years 1980–1990 were declared as the U.N. Water Decade. This influenced the work of governments and international agencies in the drive to provide clean and affordable water supplies to all. As part of this drive, international collaborative efforts led by UNICEF resulted in the sinking of millions of tubewells to extract groundwater, which is typically of superior microbiological quality when compared to surface water.
In Bangladesh alone some 10 million tubewells were drilled. UNICEF figures published in 1998 show a 47 percent reduction in infant and under 5 mortality rates in Bangladesh in the years 1980–1996 (UNICEF, 1998). However, many of the tubewells were unknowingly sunk into arsenic-contaminated aquifers. As a consequence, although microbiologically superior water supplies were obtained, some 40 million people subsequently were exposed to toxic levels of arsenic, sometimes exceeding the World Health Organization (WHO) Guideline value by a factor of 20 or more. Of 64 districts within Bangladesh, 59 were reported to contain unsafe levels of arsenic (Caussy and Priest, 2008). The situation in Bangladesh has been described as the biggest case of mass poisoning in recent history.
High levels of arsenic in drinking water have also been found elsewhere in Asia (e.g., Cambodia, China) as well as in the USA and South America. The WHO and U.S. Environmental Protection Agency (EPA) recommended limit for arsenic in drinking water is currently 10 µg/L (WHO, 2006).
There is clearly a pressing need to mitigate the problem. However, it is not so much the difficulty of removing arsenic from water, as it is the extremely low levels to which it must be reduced to ensure safety that presents the challenge to water treatment initiatives, especially in developing countries where the issues of cost and expertise often make “high-tech” solutions impractical. The challenge is not only to provide cheap and efficient treatment but also to develop technologies that are suitable for the large number of small and isolated point sources that are typical of affected rural populations.
The IUPAC Project
By the end of the millennium, the problem of arsenic contamination, particularly in Bangladesh was well documented. The implications to human health from arsenic exposure have been widely investigated by the scientific community. However, from the perspective of those living in affected areas, these problems are not fully understood and solutions for mitigation have not been adequately evaluated or communicated. It is against this background that an IUPAC project (#2003-017-2-600), established under the auspices of the Chemistry and the Environment Division, aims to accomplish the following:
- provide a global overview of natural and anthropogenic sources of arsenic in the water environment and to assess its behavior
- critically review field testing kits intended to provide cheap, quick, on-site measurements so that contamination can be accurately mapped and remediation efforts effectively monitored
- assess the health risk of arsenic contamination, with special reference to technical challenges for optimizing arsenic remediation measures that are acceptable to affected communities
- conduct critical analysis of appropriate methodologies, evaluate their suitability for different situations, and address the transferability of specific technologies that are currently associated with local conditions
For mitigation/remediation action to be taken, the interrelationships and complexities inherent to the processes need to be recognized and understood. To help achieve this, a conceptual model describing these relationships is presented (figure 1). It serves as a framework for the integration of these factors and as a basis for discussion in this article.
|Figure 1: A conceptual model describing the factors and relationship informing the decision-making process for mitigation/remediation of human arsenic exposure. (Garelick et al., 2006)
The environmental presence of arsenic derives from both natural and anthropogenic sources. Globally, most arsenic contamination of water occurs naturally when aquifers pass through bedrock containing arsenic. Significant amounts of arsenic are also introduced into the environment from anthropogenic sources, metal mining and smelting being the most important (Garelick et al., 2008).
The main regions of the world affected by arsenic contamination are depicted in figure 2, which also indicates the scale of chronic human exposure. Four principal sources of arsenic are also shown. As can be discerned from the map, areas having natural (geological) contamination of aquifers and regions with mining sources dominate. For coal burning sources, only the Guizhou region of China is shown since this region has the most serious documented problems of this type. Arsenic has also been introduced into the environment through extensive use of arsenical compounds in pesticides. It is also a constituent of wood preservatives (e.g., as copper chromate arsenate). Many sites used by the timber industry to treat lumber have been contaminated with arsenic. Agricultural and/or wood preserving sites of contamination have been reported in South Africa, Zimbabwe, and Australia, but are not included on the map as they are small localized areas.
Analysis and Environmental Sampling
Arsenic’s highly variable distribution in the environment makes it difficult to establish reliable environmental sampling regimes that accurately portray the true status of arsenic contamination. In addition, accomplishing rapid on-site sample analyses at acceptable costs remains a challenge even when reliable sampling has been achieved. A review of these methods is provided by Feldmann (2008). Monitoring for arsenic has historically been centered on drinking water, but more recently monitoring of soil and staple foods has become important because of their potential contribution to human arsenic exposure. On-site analyses of water, soil, sludge, and foodstuffs are essential to ensure that environmental risk assessments properly reflect arsenic residue burdens. Further, those undertaking environmental risk analyses have to consider additional criteria as discussed below.
Exposure, Risk, and Health Aspects
A range of arsenic species is present in the environment, which exhibit complex environmental behavior and large variations in toxicity. The pronounced differences in relative toxicities of six arsenic species are demonstrated in figure 3. The figure, although based on acute studies, further highlights significant differences in toxicity between organic and inorganic species of arsenic. These differences are also supported by a number of field studies, although it is less clear whether these differences are similarly pronounced under conditions of chronic exposure.
|Figure 3: Relative toxicity of arsenic species (MMA = monomethylarsonic acid, DMA = dimethylarsinic acid). Adapted from Vaughan (2006). [Arsine toxicity is based on rat LD50 of 3 mg/kg body weight normalized to 1].
Humans may be exposed to inorganic arsenic via air (significant intakes by inhalation occur in residents living near industrial sources where exposures to arsenic trioxide are possible), drinking water, food, and soil. However, there are strong indications that consumption of arsenic polluted water is the most important contributor to arsenic-related morbidity. The areas exhibiting the highest levels of human exposure, such as Bangladesh and West Bengal (figure 2), are known to also have the highest occurrence of arsenic contamination of ground water (Caussy and Priest, 2008).
Most food products usually contain less than 250 µg/kg of arsenic. However, seafood, including demersal fish, crustaceans, and marine algae, may contain up to ~100mg/kg of arsenic. The low levels of arsenic in plants contrast with the much higher levels in soil, ~40mg/kg, and may reflect the insolubility of many arsenic-containing minerals such as pyrite under normal soil conditions. The average U.S. daily dietary intake of arsenic is estimated to be 10 to 20µg. In Japan, however, where the seafood content of the diet is high, intakes are much larger (70 to 370µg per day) (Caussy and Priest 2008). The likely predominance of the less toxic organic species of arsenic in foodstuffs may render this source less important in morbidity. However, this is an area that has been less well researched.
Standards and Guidelines
In most countries, the maximum contaminant level (MCL), which is used as an important field monitoring criterion for total arsenic in ground- or potable-water, is set at 50 µg/L. The MCL is either a guideline or an enforceable standard, depending on the country, and is the highest level of a contaminant allowed in drinking water. However, the more stringent WHO guideline value of 10 µg/L has been set as the MCL in recent years. This has now been approved by an increasing number of countries across the world, and has been enforced in the USA since the beginning of 2006. The MCL does not distinguish arsenic species, although as shown in figure 3, speciation is clearly an important factor in arsenic toxicity.
Arsenic levels in soil, sludge, and foodstuffs are less well regulated. Some countries, such as the UK and Canada, have environmental quality criteria for soil (10–50 mg/kg) that are use-dependent. Guidelines for arsenic in foodstuff have not yet been so widely agreed upon. At present, only Australia has established a guideline (1 mg/kg) for total arsenic (no discrimination among species) in foodstuff. China has so far introduced maximum levels for different foodstuffs and was the first country to incorporate arsenic speciation into its guidelines (e.g., 0.15 mg/kg inorganic arsenic for rice). However, the importance of distinguishing among arsenic species in guidelines is an area that requires further progress (Feldmann, 2008).
The health outcomes from exposure to arsenic depend not only on the source and chemical species of arsenic, but on dose, modality, and duration of exposure. Both acute and chronic intakes result in toxicity, although drinking contaminated water is only likely to produce effects under conditions of chronic intake. These effects include skin lesions, disturbances of the peripheral nervous system, anemia, and leucopenia, liver damage, circulatory disease, and cancer. Many of these effects have been observed in populations in Taiwan, Argentina, and Bangladesh that consumed contaminated water. The bioavailability of arsenic is a key determinant of the health outcomes of exposure because it is related to the ability of arsenic to be liberated from ingested matrices (e.g., soil, water and food) and thus enter into the blood stream where it exerts its toxic effects. Studies from human volunteer and animal experiments show about 90 percent of ingested, dissolved, inorganic arsenic salts are absorbed from the gut and enter the blood stream. This percentage is much higher than for most other ostensibly non-essential elements and, in part, is a feature of the similar chemical properties of the arsenate (AsO43-) and phosphate (PO43-) ions, the latter being an essential body component that is avidly absorbed from the gut (Caussy and Priest, 2008).
Substantial efforts have been invested in developing techniques for removing arsenic from water in both field and laboratory conditions. For any technology to be effective and appropriate for use in affected areas of developing countries, it should use local resources and be low-cost, versatile, transferable, and simple to operate and manage. It is essential that such technologies are accessible to local communities and especially to women. Throughout the developing world, women collect and carry water for their families, use water for cooking and cleaning and for growing food, and therefore should be at the forefront as users of arsenic treatment technologies. Even treatment systems that are highly successful from a technological viewpoint may not succeed in rural areas, such as villages of India and Bangladesh, unless they mesh with local circumstances and are well accepted by the local population.
|Figure 4. Color labelling of safe (green) and unsafe (red) tubewells in Bangladesh. [Photographs courtesy of Feroze Ahmed, Department of Civil Engineering, Bangladesh University of Engineering & Technology, Dhaka, Bangladesh.]
Although water treatment is central to arsenic pollution mitigation, it is important to recognize that a change in the source of water, rather than treatment of the water locally, may be the most appropriate solution in some cases.
Where municipal piped water supplies exist, connecting to these may provide the appropriate solution for those using contaminated private sources. Even in areas where no municipal piped supplies are provided, switching to a different source well or to harvesting rain water may still be the best solution (Jones et al., 2008). An example of how this may be managed is illustrated by the practice in rural Bangladesh of color-labelling tubewells that provide water deemed unsafe for consumption (figure 4).
Conventional treatment plants may employ several methods for removing arsenic from water. Commonly used processes include oxidation and sedimentation, coagulation and filtration, lime treatment, adsorption onto sorptive media, ion exchange, and membrane filtration. However, in the most affected regions large conventional treatment plants may not be appropriate and factors such as cost and acceptability as well as performance must be considered. A summary of mitigation options, arsenic removal performance, and relative costs are presented in table 1.
|Table 1. Evaluation of Performance and Costs of Arsenic Pollution Mitigation Options. Adapted from Visoottiviseth and Ahmed (2008) and Ellis and Garelick (2008).
Multi Criteria Analysis
The identification and selection of best practices for arsenic mitigation options requires careful evaluation of economic, technological, and environmental factors. Regional social values and inter-disciplinary and inter-institutional contexts are also important considerations. Procedures designed to resolve arsenic remediation must also address complex problems in a manner that is comprehensible to and auditable by stakeholders. Where divergent-minded stakeholders are involved, multi-criteria analysis (MCA) is an approach that can facilitate optimal decision-making.
MCA is a decision-making tool originally developed to help evaluate options that rely on values that are not easily assignable. Cost-benefit analysis is an example where a range of defining criteria (or influencing factors) important to stakeholders are considered. The facilitation of discussions and increased transparency of decision-making processes are possible through MCA since it enables stakeholders to process large amounts of complex information, thereby helping them see the “bigger picture.”
MCA procedures utilize a performance matrix (consequence table) in which options for arsenic remediation are presented (and described) in columns, and performance values for options are presented in rows against predefined criteria. Data may be entered in the matrix as utility scores in various forms, such as numerical, semi-quantitative, and/or descriptive. The performance of each option can then be aggregated against all the criteria from which an eventual overall comparison and decision on the best mitigation option can be made.
Source-exposure vector, health risk, cost, social acceptance, and technical competency are the primary MCA criteria on which a judgement or decision is reached in relation to arsenic treatment options. An example of such a matrix is shown in table 2. The table tracks relevant and appropriate properties for each of the listed criteria that may influence outcomes. Data in the performance matrix is converted into numerical values through application of a specific utility scale scoring and weighting technique for each criterion. A detailed account of these procedures is provided by Ellis and Garelick (2008).
Utility scores and/or weights will clearly influence performance matrix outcomes. Sensitivity analysis can also be performed by reiterating the analysis using different scores or weightings. This is a useful and important approach that enables the MCA methodology to be used as a negotiating tool in the decision-making process, allowing areas of stakeholder agreement and disagreement to be highlighted.
A substantial knowledge pool exists about human exposure to arsenic contamination; the scale of contamination in many regions is well documented, human health effects are known and understood and potential technological solutions have been elucidated. However, the implementation of successful mitigation requires this knowledge to be applied more effectively. This can only be achieved by enhancing the understanding of the interactions between socio-economic and scientific issues, thus enabling practitioners to develop effective and appropriate mitigation programs.
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Garelick, H. et al. (2008). “Arsenic Pollution Sources.” Reviews of Environmental Contamination and Toxicology. 197. [In Press].
Jones, H. et al. (2008). “Case Report: Arsenic Pollution in Thailand, Bangladesh and Hungary.” Reviews of Environmental Contamination and Toxicology. 197. [In Press].
Keusch, G.T. et al. (2006). “Diarrheal Diseases.” In Jamison et al (eds): Disease Control Priorities in Developing Countries (2nd Edition), ed. , 371–388. New York: Oxford University Press. <www.dcp2.org/pubs/DCP>; Chapter 19.
UNICEF (1998) The State of the World’s Children 1998.
Vaughan, D.J. (2006) “Arsenic” Elements 2: 71–75.
Visoottiviseth, P. and Ahmed, F. (2008) “Technology for Remediation and Disposal of Arsenic.” Reviews of Environmental Contamination and Toxicology. 197. [In Press].
WHO (2006) Guidelines for Drinking-Water Quality. Volume 1—Recommendations. First addendum to third addition. <www.who.int/water_sanitation_health/dwq/gdwq0506.pdf>
WHO (2008) “Water Quality Interventions to Prevent Diarrhoea: Cost and Cost-Effectiveness” (T.F. Clasen, L. Haller) Geneva 2008 WHO/HSE/WSH/08.02
*This article arose from work carried out by an IUPAC project (2003-017-2-600) sponsored by the Chemistry and the Environment Division (VI). The participants were F. Ahmed, G. Borbély, M.K. Bux, D.H. Caussy, A. Dybowska, J.B. Ellis, J. Feldmann, R. Földényi, Z. Galbács, H. Garelick, H. Jones, N. Kováts, N. Priest, Y. Shevah, E. Valsami-Jones, and P. Visoottiviseth.
Hemda Garelick <H.Garelick@mdx.ac.uk> is a professor in the School of Health and Social Sciences at Middlesex University, The Burroughs, London, United Kingdom. Huw Jones <H.Jones@mdx.ac.uk> is a professor in the Department of Natural Sciences at Middlesex University.
2. Worldwide distribution of arsenic contaminated regions,
showing source of arsenic and numbers of people at risk
of chronic exposure. Note: estimates of people at risk
of poisoning are very difficult to quantify, particularly
in areas where geochemical surveys are limited. These
estimates are broad and are based on four criteria: 1)
prevalence of current recorded cases of arsenicosis, 2)
likelihood of ingested concentrations exceeding 50 µg/L)
number of people living in exposed areas, 4) likely ability
of region to mitigate/remediate against contamination.
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8. Webster JG & Nordstrom DK (2003). Geothermal arsenic. In: Welch AH, Stollenwerk KG, (Eds) Arsenic in Ground Water: Geochemistry and Occurrence, Kluwer Academic Publishers Boston, pp101-126.
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10. Wilson FH & Hawkins DB (1978). Arsenic in streams stream sediments and ground water, Fairbanks area, Alaska. Environ Geol 2:195-202.
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13. Thornton I & Farago ME (1997). The geochemistry of arsenic. In: Abernathy CO, Calderon RL, Chappell WR (Eds) Arsenic Exposure and Health Effects. Chapman & Hall London pp 1-16.
14. Kurttio P, Pukkala E, Kahelin H, Auvinen A & Pekkanen J. (1999). Arsenic concentrations in well water and risk of bladder and kidney cancer in Finland. Environ Health Perspect 107(9):705-710.
15. Filippi M, Golias V & Pertold Z (2004). Arsenic in contaminated soils and anthropogenic deposits at the Mokrsko Roudny´ and Kas¡perske´ Hory gold deposits Bohemian Massif (CZ). Environ Geol 45:716-730.
16. Földényi, R , Kováts, N, Borbély, G & Galbács, Z (2008). Arsenic in drinking water supply in Hungary. Rev. Environ. Contam. Toxicol. 197 (in print)
17. Gurzau ES & Gurzay AE. (2001). Arsenic in drinking water from groundwater in Transylvania, Romania 181-184. In Chappell WR Abernathy CO, Calderon RL, (Eds) Arsenic Exposure and Health Effects IV. Elsevier, Amsterdam.
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19. Modabberi S & Moore F (2004). Environmental geochemistry of Zarshuran Au-As deposit NW Iran. Envirol Geol 46:796-807.
20. Shrestha B. (2002). Drinking water quality: future directions for UNICEF in Pakistan. Consultancy report 2 of 3, Water Quality, SWEET Project, UNICEF Pakistan, Islamabad.
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34. Smith E, Smith J, Smith L, Biswas T, Correll R & Naidu R. (2003). Arsenic in Australian environment: an overview. J Environl Sci Health, Part A: Toxic/Hazardous Substances and Environmental Engineering, A38 (1): 223-240.
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|Table 2 provides an example of the MCA procedure in which a performance matrix is constructed from
utility scores for each key performance indicator. [View full-page pdf.]
Page last modified 5 August 2008.
Copyright © 2003-2008 International Union of Pure and Applied Chemistry.
Questions regarding the website, please contact firstname.lastname@example.org |
Garbl's Writing Center
cache, cachet Sometimes misused or confused nouns. A cache (pronounced "cash") is "a hidden supply of weapons, valuables and other things"; "a secret place for hiding things"; and "memory for computer data." A cachet (pronounced "ca-shay") is "a quality that people admire," "an official seal," and "a commemorative mark."
CAD An acronym for computer-aided design. Spell out on first reference.
calendar Commonly misspelled.
call letters Use all caps when referring to TV and radio stations and networks. Use hyphens to separate the type of station or network from the basic call letters: KMTT-FM, KIRO-AM, KING-TV, CBS-TV. Don't italize call letters or put them in quotation marks. Use these formats for other types of stations that mix numbers and letters: WXY12, N3OPQ. See channel, citizens band, station.
callous, callus Sometimes confused. Callous is an adjective meaning "hardened and uncaring of others' suffering." Callus is a noun for "an area of hard, thick skin."
calvary, cavalry Often confused. Capitalized, Calvary is "the place where Jesus was crucified." Lowercased, calvary is "an outdoor depiction of the crucifixion" and "an experience involving intense suffering." Cavalry is a noun for "military troops on horseback or in armored vehicles."
camaraderie Commonly misspelled. Don't confuse with spelling of comrade.
cancel, canceled, canceling, cancellation Commonly misspelled.
can, could Use can to express certainty or willingness in being able to do something. Use could when there's less certainty or when doing something depends on something else. See can, may below; may, might; will, would.
can, may Commonly confused. Use can when writing about capability, physical or mental ability, or the power to do something. Use may when writing about authorization or permission and sometimes possibility: They can finish the report by November. May we have an extra month to finish the report? You may lead the horse to water, but you can't make it drink. May is almost always the correct word to use in a question. See can, could above; may, might.
cannon, canon Often confused. A cannon is a "a large, mounted piece of artillery; a big gun." A canon is "a law of a church," "an accepted rule or principle of behavior," and "a set of literary works."
cannot One word.
canvas, canvass Sometimes confused or misspelled. A canvas is "a heavy, closely woven, coarse cloth used for tents, sails, bags, oil paintings and other things." To canvass is "to examine or discuss something (like votes) in detail," "to go through an area asking people for votes or opinions," and "to sell something house to house."
capacity See ability, capability, capacity.
capital, capitol Often confused or misspelled. Capital is a city, the seat of government. Do not capitalize: Salem is the capital of Oregon. Capital city is redundant. Capital also refers to money. Capitol is the building in which the U.S. Congress or the state Legislature meets. Capitalize U.S. Capitol and the Capitol when writing about the building in Washington, D.C., and do the same when writing about state capitols: The California Capitol is in Sacramento. Capitol building is redundant. See Capitol Hill.
capitalization Rule No. 1: Use capital letters to begin proper nouns, sentences, headings, some abbreviations and acronyms, and the important words in composition titles. Proper nouns are the particular names of people, places and things. Rule No. 2: Do not capitalize the first letter of a word (or words in a phrase) simply to highlight it or because you or someone else think it's an important word. Excessive, arbitrary capitalization distracts the reader and hinders reading.
Check this or another style manual for capitalization of a particular word or type of word. If not listed there, check your dictionary. And if still in doubt, lowercase.
Except for acronyms and some abbreviations, avoid capitalizing all the letters in a word, sentence, heading, headline or phrase--including brand names, logos and trademarks. For emphasis, try other typographical uses instead, such as boldfacing, italics, color, type size and different but complementary typefaces. Also see headlines, headings; underlining.
Capitalization of abbreviations and acronyms varies. For guidance, see abbreviations and acronyms, entries in this style guide for specific words and terms, and your dictionary. Although the abbreviation or acronym is capitalized for some common or generic nouns and terms, lowercase the spelled-out form; for example, see environmental impact statement.
Capitalize the first word of every sentence, heading and headline, including quoted statements and direct questions. Even if a person, business or organization begins its name with a lowercase letter, capitalize the first letter of the name at the beginning of sentences, headings and headlines: Gary de Shazo won the design award. De Shazo expressed appreciation for the support of his colleagues. Also see composition titles.
Capitalize proper nouns that specifically name a person, place or thing, unless a person, business or organization requests a lowercase first letter. If a name begins with a lowercase letter, capitalize the first letter of the name at the beginning of sentences and headlines.
Capitalize common nouns such as party, river and street when they are an integral part of the full name for a person, place or thing: Ballinger Street, Rheinard River, Queens County, Democratic Party, Puget Sound. Lowercase those common nouns when they stand alone in later references: the party, the river, the county, the street, the sound.
Lowercase common noun elements of names in all plural uses: Democratic and Republican parties, Ackley and Messer streets, 154th and 156th avenues southeast. But don't lowercase the common nouns when the form is not plural: Your sister can catch a bus on First or Third Avenue.
Capitalize the proper names of nationalities, peoples, races, tribes and so on: African American, American Indians, Arab, Asian, Jewish, Latino, Muckleshoot, Tulalip, Puyallup. Lowercase black, white, red and so on. See race.
Organizations should adopt specific capitalization guidelines for their governing boards, facilities, job titles and descriptions, organizational structure, and programs, projects and plans. It's efficient to develop styles consistent with a standard, readily available, published reference source. For recommended capitalization guidelines, check individual items in this style manual or see the items below: committees, facilities, job titles and descriptions, organizational structure, and programs, projects and plans.
carat, caret, carrot, karat Often confused or misspelled. A carat is "a unit for measuring the weight of jewels," while a karat is "the unit for measuring the purity of gold." A caret is "a mark [like this > symbol but pointing up] used in editing and proofreading to show where something is to be inserted." A carrot is "a long, thick, orange vegetable."
cardiopulmonary resuscitation See CPR.
careen, career Sometimes confused verbs. To careen is "to lean or tilt sideways, especially a ship in high winds or on a beach for cleaning and repairs" and "to sway or lurch from side to side." To career is "to move at full speed ahead, to rush wildly."
carpool One word. It may be used as a noun, verb or adjective: Her neighbors formed a carpool to save gas and money. They carpooled to work to save gas and money. She requested some carpool information. See high-occupancy, vanpool.
cast The past tense of this verb is also cast, not casted.
catalog, catalogue Both are correct, but catalog is commonly preferred.
catch-22 Sometimes misused. A catch-22 is not any simple catch, or any tricky situation with a hidden complication. From the excellent antiwar novel of the same name by Joseph Heller, a catch-22 is an absurd or paradoxical situation in which the desired outcome is impossible because of built-in illogical rules: The experienced editor couldn't get promoted to supervisor because he didn't have any experience as a supervisor.
catchup, catsup See ketchup.
category Overstated. Simplify. Try replacing with group.
cavalry See calvary, cavalry.
CB See citizens band.
CBD See central business district.
CD-ROM Acronym for compact disc read-only memory. The acronym is acceptable on first use. CD-ROM disc is redundant.
cease, seize Sometimes misused. To cease is "to bring something to an end." To seize is to "take hold of something or someone suddenly and forcibly," "capture a place or assume control using force," "confiscate," or "take initiative eagerly." The saying is seize the day, not cease the day. Also, cease is overstated and formal. Simplify. Try stop, end or finish instead.
ceiling Sometimes misused. Before using this word, look up. You're looking at the ceiling. As the upper limit on something, ceiling is also a useful metaphor for maximum or limit: a ceiling on taxi rates. You can raise a ceiling, lower it or even remove it. But you could mangle a ceiling if you try increasing it, decreasing it or waiving it. See target.
cellphone, cellular phone, cellular telephone Cellphone (one word) is acceptable on first reference. Also, smartphone is one word.
Celsius Use this term instead of centigrade for the temperature scale that is part of the metric system. Spell out and capitalize on first reference. The abbreviation C (capitalized, no period) may be used on second reference with a numeral: The temperature dropped to 5 C Monday night. See Fahrenheit, temperatures.
cement, concrete Often confused. Cement is dry, powdery ingredient of concrete. Concrete--a mixture of cement, water and sand or gravel--is used to form pavement, blocks, walls, driveways, sidewalks and so on.
cemetery Commonly misspelled. Memory aid: Almost every other letter is an e.
center around Illogical and redundant. Substitute on, in or at for around, or use revolve around. Avoid center upon.
center stage Cliche. Omit, or try prominent, center of attention or focus of interest instead.
cents Spell out and lowercase cents using figures for amounts less than a dollar. Use the $ sign and decimal system for larger amounts: 33 cents, $2.04, $3.47. Do not use zeros if there are no cents: $8, not $8.00. Avoid using the cent symbol: ¢. But if you must use it, be careful not to use the redundant .33¢ or $.33¢. See dollars.
CEO Abbreviation for chief executive officer. Acceptable on first use before a name or standing alone, if spelled out somewhere in a document. Spell out less familiar chief operating officer and chief financial officer.
chair, chairman, chairperson, chairwoman Use chair as the title for the heads of councils and committees, unless the person in the position prefers chairman, chairwoman or chairperson. Capitalize as a formal title before a name. Do not capitalize as a casual, temporary position.
chaise longue Sometimes misspelled (and mispronounced). It's French for "long chair," with a back support and seat long enough to support outstretched legs. Don't spell it chaise lounge (And don't pronounce it "chase lounge." Say "shayz long.")
changeable Commonly misspelled
changeover One word, no hyphen.
chapter Capitalize when used with a number to name a section of a book or legal code: Chapter 11. Lowercase when standing alone.
character Commonly misspelled.
charts, tables Charts and tables are useful in reports to present information concisely. Abbreviations not typically used in text are acceptable in charts and graphs because of limited space. But abbreviations must still be clear to the reader and consistently used. Also, charts and graphs should have titles. Capitalize the first letter of proper nouns and key words in the titles and headings of charts and tables. Type styles and formats used in charts should be consistent throughout a publication. When using several charts or tables, assign numbers. When mentioning a chart or table in the text, capitalize the word chart or table and use the numeral: As Table 6 shows, traffic congestion has gotten worse since they built the football stadium.
chat room Two words.
chauvinism Sometimes misused or confused. It used to mean only "excessive, unreasoning, or blind devotion to one's country, a fanatical patriotism." But it's now commonly applied to excessive pride in a person's group, race or sex, especially males. If you must use this derogatory word, be clear about whom or what you're describing.
check in (v.), check-in (n. and adj.)
check out (v.), checkout (n. and adj.)
check up (v.), checkup (n.)
Chicano See Hispanic, Latino.
chief Capitalize as an official job title before a name: Facility Maintenance Chief Suzanne Zentin. Lowercase when used alone or after a name between commas: She called Leif Elliott, customer services chief, about the complaint. See titles.
chief executive officer See CEO.
child care, child-care Hyphenate as an adjective: He uses a child-care agency in downtown Olympia. Don't hyphenate as a noun: He searched everywhere for the best child care.
childish, childlike Sometimes confused or misused adjectives for describing behavior typical of children. Use childish to describe unfavorable qualities like immature and silly. Use childlike to describe favorable qualities like sweet, innocent and trusting.
children Usually, use first names on second reference for children 15 or younger. For older children, the last name is usually suitable. Although it may not assure mature behavior, treat people 18 and older as adults; use their last names on second reference.
children's The apostrophe always goes before the s when showing the possessive: the Children's Home Society. Don't use childrens' (with the apostrophe after the s); children is already plural.
choice between, choose between When between follows choice or choose, use and, not or, between the choices: The students had a choice between taking a midterm exam and finishing another homework assignment. We had to choose between a helicopter ride and a catamaran ride.
chord, cord Often misused nouns. A chord is "a combination of three or more musical notes played at the same time." A cord is "a piece of wire covered with plastic for carrying electricity," "a measure of wood cut for fuel," "a ribbed cloth," "a piece of thick string or thin rope," and "a part of the anatomy resembling a cord": vocal cords.
Christmas See holidays.
cities and towns Capitalize the names of cities and towns in all uses. Capitalize city as part of a proper name: New York City, Kansas City.
Lowercase city when used as an adjective or noun: the city budget, mayor of the city. Capitalize city when mentioning the proper name of a governmental unit: He worked for the City of Kennewick. But lowercase city--or omit the redundant city of--when naming cities in other uses: They visited the city of Edmonds. They visited Edmonds.
Lowercase general descriptions such as north Seattle. Capitalize widely recognized names for the sections of a city: Laurelhurst, Magnolia, West Seattle, Rainier Beach and the University District.
citizen A citizen is a person who has the full civil rights of a nation through birth or naturalization. Cities and states in the United States do not grant citizenship. Use resident to include noncitizens as inhabitants of states, cities and communities. See public.
city council Capitalize when part of a proper name: The Langley City Council scheduled a meeting. Also capitalize if the name of the city is clear: The City Council passed a motion. Lowercase in other uses: the council, the Langley and Coupeville city councils.
citywide One word.
class See collective nouns.
clean bill of health Cliche. Try good report, good condition, doing well, fine, healthy or strong instead.
clean up (v.), cleanup (n. and adj.) The cleanup lasted three months. It took three months to clean up the river. Also, think about dropping up from clean up.
cliche William Safire, Fumblerules, 1990: "Last but not least, avoid cliches like the plague." And if you must use a cliche, don't put quotation marks around it.
climatic, climactic Occasionally confused adjectives. Use climatic when describing the climate or changes in the weather. Use climactic when describing a climax, key dramatic moment or highest point.
climb down, climb up Climb up is usually redundant, and climb down seems illogical. So use climb alone to mean "going up" and climb down as an acceptable idiom for "going down." Both terms are preferable to the formal ascend and descend.
close proximity Redundant and wordy. Simplify. Use near or close to instead. Also, closeness and nearness are both preferred to the formal proximity.
clothes Sometimes misspelled as the soundalike close. Memory aid: Clothes are made of cloth, which wears a th.
co- Hyphenate when forming nouns, adjectives or verbs that show occupation or status: co-host, co-pilot, co-signer, co-worker. Omit the hyphen in other combinations, including coordinate, coordination, cooperate, cooperation, cooperative. See prefixes.
coast Lowercase when writing about the physical shoreline: Atlantic coast, Pacific coast. High winds battered the Atlantic coast. Capitalize when writing about regions of the United States lying along such shorelines: The Atlantic Coast states all supported the Democratic candidate. The Pacific Coast states all had sunny weather.
Do not capitalize when writing about smaller regions: She loves the Oregon coast in November.
coed Don't use the outdated, sexist coed (or co-ed) as a noun to refer to a female student. Coed and coeducational are fine as adjectives to note that both sexes are involved.
coliform bacteria Bacteria common to the intestinal tract of people, other mammals and soil. Always lowercase.
collaborate, corroborate Occasionally confused verbs. To collaborate means "to work together for a special purpose" and "to cooperate with an enemy." To corroborate means "to confirm one statement by referring to another statement."
collectible This spelling more often preferred than collectable. Both spellings carry the same meaning.
collective nouns Collective nouns name a group or collection of people, places, things, ideas, actions or qualities, including board, class, committee, crowd, family, group, herd, jury, panel, public, orchestra, staff, team. Nouns that show a unit take singular verbs and pronouns: The board is electing its committee chairs. The crowd is eager to march. To stress individuals in a group, use members of: Staff members answered questions. Some members of the panel left before lunch. See it.
Some nouns are both singular and plural in meaning, including corps, chassis, deer, fat, fish, grease, moose, oil, public, sediment, sheep, soil, water and waste. The use of a singular or plural verb in a particular sentence conveys the meaning. Because these words are already plural, avoid adding s or es to make them plural: Scientists studied sediment from Charger Bay. The geologist took samples of soil from the site. When mentioning various types or species, however, plural spellings may be used: Scientists studied Fox Lake and Lake Roosevelt sediments. The site contained both glacial and sandy soils.
Follow the rules of subject-verb agreement when using the proper names of athletic teams and musical bands or groups: The Seattle Mariners are on the road. The Seattle Storm is an event sponsor. The Beatles were wonderful at the old Seattle Center Coliseum and so were the Rolling Stones. The Who is still terrific.
college names See university names.
collide, collision Two objects must be moving before they can collide. An accident involving a moving car and a stationary telephone pole is not a collision, for example; it's a crash. See near miss, near-miss.
colon (:) The colon has three main uses, all of which involve pointing the reader toward the words that follow the colon. The colon always follows a whole sentence in these uses. Don't combine a dash and a colon.
The most frequent use is to introduce a list, often after expressions such as the following or as follows: Loretta Schwieterman appointed three people to the committee: David Allen, Greg Edwards and Jean Rheinhard. The Parks Department has scheduled open houses in the following communities: Valley View, April 5; Gantry, May 6; and Sierra Hills, Aug. 7.
Don't use a colon immediately after a verb. Incorrect: Loretta Schwieterman appointed: David Allen, Greg Edwards and Jean Rheinhard to the committee. Correct: Loretta Schwieterman appointed David Allen, Greg Edwards and Jean Rheinhard to the committee. For more information on creating lists, see lists, semicolon.
Second, the colon can be used to stress the word, words or sentence that follows it: He had only one thing on his mind: flowers. The news was good: No one would be laid off. When used this way, the colon replaces such words as that is, namely and for example. Capitalize the first word after a colon if it is a proper noun or the start of a whole sentence.
Third, use a colon to introduce a quotation longer than one sentence within a paragraph and to end a paragraph that introduces a quotation in the next paragraph. Use a comma, however, to introduce a quotation of one sentence that stays within a paragraph. See attribution, comma below, quotations, quotation marks.
color Usually redundant and wordy when naming a color. Simplify. Try dropping in color, colored and the color from phrases like blue in color, red colored, the color green.
combine together Redundant and wordy. Simplify. Drop together.
coming Often misspelled. It has one letter m, not two.
comma (,) The following guidelines treat frequent questions about eight essential uses of the comma.
First, in a series of three or more terms with a single conjunction, use a comma after each term: She opened the closet, grabbed a coat, and picked up an umbrella. In a complex series of phrases, the serial comma before the final conjunction aids readability. In a simple series, the comma is optional before the conjunction: The van is economical, roomy and dependable. Also, put a comma before the final conjunction in a series if an integral element of the series needs a conjunction: He likes folk, rock, and rhythm and blues. Don't put a comma before the first item in a series or after the and in a series. See lists, semicolon.
Second, use a comma to join two independent clauses with a conjunction. An independent clause is a group of words that could stand on its own as a complete sentence; it begins with its own subject. The most common conjunctions are but, and, for, nor, or, so and yet: The council's Water Resources Committee will go over the resolution Jan. 12, and the full council is scheduled to act Feb. 11. Don't create run-on sentences by combining two or more independent clauses with only commas. Either insert conjunctions after the commas or break the clauses into separate sentences. See sentence length.
Third, use a comma to separate an introductory phrase or clause from the rest of the sentence: After graduating from college, he joined AmeriCorps. It may be omitted after short introductory phrases (less than three words) if no ambiguity would result: On Thursday the Kennewick City Council will decide the issue. When in doubt, use the comma, especially when it separates two capitalized words.
Fourth, enclose parenthetic expressions between commas. Parenthetic expressions are word groups that are not essential to the meaning of a sentence. If a parenthetic expression is removed, the sentence would still make sense: The social services manager, who toured the Snoqualmie Valley last week, will make her recommendations today. They took one of their sons, Leif, to the concert. His wife, Donna, is a middle school teacher. As shown in the examples, commas always go both before and after a parenthetic expression within a sentence. If you'd prefer to stress a parenthetic phrase, put it between dashes; you can play down such a phrase by placing it between parentheses. Also see this, that, who whom.
Also use commas to set off a person's hometown when it follows the name: Rachel Solomon, Danbury, opened a new restaurant. If using a person's age, set it off by commas: Tom O'Rourke, 69, opened a new restaurant.
Do not use commas to set off an essential word or phrase from the rest of a sentence. Essential words and phrases are important to the meaning of a sentence: They took their daughter Jennifer to school. Their son Nils works at Ticketmaster. (They have more than one daughter and more than one son.)
Fifth, use commas to set off words and phrases such as however, meanwhile, in fact, in addition, moreover, nevertheless, as a result, thus, therefore, for example, finally and in other words. Usually, place a comma after such expressions when they begin a sentence, and place commas before and after the expressions when they are within a sentence. See however, in fact, in addition to, moreover, nevertheless.
Sixth, use commas to separate a series of adjectives equal in rank. If the adjectives could be rearranged without changing the meaning of a sentence or if the word and could replace the commas without changing the sense, the adjectives are equal: A sleek, new car. A thick, black cloud. See hyphen.
Use no comma when the last adjective before a noun outranks its predecessors because it is an integral element of a noun phrase: a silver articulated bus.
Seventh, use a comma to set off a direct one-sentence quotation within a paragraph: Theodore Roosevelt said, "It's not the critic who counts." Use a comma before the second quotation mark in a quotation followed by attribution: "No comment," said Jerry Carson. See attribution, punctuation, quotations, quotation marks.
And eighth, use a comma to separate the parts of numbers, dates and addresses. Use a comma for figures higher than 999: More than 5,000 people attended the event.
Use commas to set off the year in complete dates: The department released its report Nov. 16, 2002, for public review. But don't separate the month from the year when not using a date. They held their first retreat in January 1994.See dates.
Use commas to set off cities from names of states or nations: She went to Vancouver, Wash., to tour the bridge retrofit program. He traveled to Paris, France, on vacation.
commence See begin, commence, start.
commitment Commonly misspelled. Remember the root word, commit, with two m's and one t.
committee Commonly misspelled. Capitalize if part of the proper name: the Langley City Council's Human Services Committee. Lowercase when used alone: The committee passed the motion. See capitalization, collective nouns, subcommittee.
common, mutual They have a subtle difference in meaning. Use common to describe something shared by two or more people or things: a common goal, common interests. Use mutual to describe a feeling or action that's exchanged or reciprocal between two or more people or things: mutual respect, mutual efforts. See mutual.
community action grant Lowercase. Avoid abbreviating CAG.
compact disc CD is acceptable on later references.
company names When using a company (or product) name, you have no obligation to help a company market itself (or its products). For most proper names, capitalize the first letter of each word, or capitalize a different letter if preferred by a company: eBay. But capitalize the first letter if it begins a sentence. Do not use all capital letters unless the letters are individually pronounced: IBM and BMW but Subway and Ikea (not SUBWAY and IKEA). Don't use exclamation points, asterisks and plus signs that some companies use in logos and marketing materials for their company (and product) names: Yahoo, not Yahoo!; Toys R Us, not Toys "R" Us. Unless it's part of a company's formal name, replace the ampersand (&) with and.
Abbreviate company, corporation, incorporated and limited when using them after the name of a corporate entity: the Boeing Co., American Broadcasting Cos., Chevron Corp. Don't use a comma before Inc. or Ltd. even if it's included in the formal name. Do not abbreviate those words in business correspondence. In business correspondence, spell out those words when part of the proper name: the Boeing Company. See firm, incorporated.
If company, companies or corporation appears alone in second reference, spell out and lowercase the word: The company showed a loss in the third quarter.
The forms for possessives: the Boeing Co.'s profits, American Broadcasting Cos.' profits, Chevron Corps.' profits.
comparable Commonly misspelled.
compare and contrast Probably one of your first school lessons in writing redundantly involved essays to compare and contrast things. To compare things is to discover and describe their similarities and differences. You don't also have to contrast them.
compared with, compared to Often confused. The more common phrase, compared with means "to examine the similarities or differences of two or more things": He averaged 23 points a game in 2001 compared with 17 points a game last year. The speaker compared Congress with the British Parliament. The less common compared to means "to liken two or more things, say they are similar or show a resemblance": The backhoe operator compared her work to climbing Mount Everest. He compared life to a battle. Memory tip: Compared to is metaphorical while compared with is statistical.
compass directions See directions and regions.
compatible Commonly misspelled. Remember the relationship is not able, it's ible.
compensate Unless you're paid by the syllable or letter, simplify and use pay, if that's what you mean.
complacent, complaisant Sometimes confused adjectives. Use complacent to describe someone who's satisfied and content with his or her accomplishments. Use complaisant to describe someone who's willing and eager to please.
complement, compliment Often misused or confused. Complement is a noun or verb for "something that fills up or completes": The company has a complement of 250 drivers, 75 mechanics and 10 office workers. The two ideas complement each other well. A hat may complement a suit, but you would compliment the wearer on her or his hat. A related term: full complement.
Compliment is a noun or verb for "praise or a flattering remark" and "something free": The supervisor complimented the staff for a job well done. The supervisor's compliment boosted morale.
complete (v.) Overstated and formal. Simplify. Try replacing with end or finish unless you're writing about filling in missing or defective parts. Or try replacing with fill in or fill out.
completely This adverb is often completely redundant. Simplify. Don't use completely before full and words like dedicated, destroy, devoted, eliminate, perfect, silent, superfluous, unanimous and unique--and redundant.
comply with Try replacing with simpler follow, keep to, meet or obey.
component Overstated. Simplify. Change to part or ingredient.
compose, comprise, include Compose is not synonymous with comprise. Compose means to create or put together: The division is composed of six sections. Compose takes of, but comprise never does.
Comprise means to contain, consist of or embrace. The whole comprises the parts. Use it in the active voice and name all the parts that make up the whole after the verb: The division comprises six sections. The zoo comprises mammals, reptiles and birds. Don't use comprised of. Think about using simpler consist(s) of or contain(s).
Use include when what follows is only part of the whole: city government includes the Parks and Human Services departments. See constitute.
composition titles Capitalize the first letter of main words in titles of books, long poems, long musical compositions, magazines, movies, newsletters, newspapers, plays and works of art such as paintings and sculpture. Italicize the names of such works, or underline them if italic type is not available.
Use a colon between a book's title and its subtitle: Woe is I: The Grammarphobe's Guide to Better English in Plain English.
Capitalize the first letter of main words and enclose in quotation marks the titles of dissertations, essays, lectures, short musical compositions, short poems, short stories, songs, speeches, radio and television programs, articles in periodicals and chapters of books. If the title is part of a sentence, commas and periods go inside the closing quotation mark. Other punctuation, such as the question mark and the exclamation point, goes inside the quotation mark if it's part of the title; if it applies to the entire sentence, it goes outside the quotation mark.
Capitalize--but don't italicize, underline or enclose in quotation marks--the names of brochures, bulletins, forms, reports, software, websites, and catalogs of reference material, such as almanacs, directories, dictionaries, encyclopedias, gazetteers and handbooks.
When capitalizing hyphenated words in a title, choose a style and follow it consistently. Simplest is to capitalize only the first word unless later words are proper nouns or adjectives: Unique benefits for part-time violinists, All-American flag-waving techniques. Second is to capitalize all words except articles, short prepositions and short conjunctions: Over-the-Counter Acid Reducers for Sale Here, A Matter-of-Fact Approach to Guitar Tuning, A New Park-and-Ride Lot for Commuters. Optional exceptions to the second style are to lowercase the word after a prefix unless it is a proper noun or adjective and to lowercase the second word in a spelled out number: Anti-intellectual Conduct, Twenty-first Century Values.
compound words Compound words are formed differently for different parts of speech. When forming a compound, such as start up or start-up, first determine the part of speech you want, such as a noun, adjective or verb. Then check your dictionary and style manual for the correct spelling. If not listed in either source, follow these guidelines (also see hyphen, initial-based terms):
conceal Overstated and formal. Simplify. Try replacing with hide.
concept Overstated. Simplify. Change to idea or design.
concerning Overstated and formal. Try replacing with about.
concise (adj.), concisely (adv.), conciseness (n.) Means "brief and to the point, short and clear." To be concise, write only what you must to make your point, removing all unnecessary words and details. Succinct writing is clear and precise using the fewest words possible. Pithy writing is compact but also meaningful and witty. Terse writing is concise and polished but potentially curt in its brevity. Laconic writing is brief but also rude or ambiguous, mysterious and uncommunicative.
Writing that is not concise may be wordy (more words than necessary) and verbose (obscure and tedious), rambling (aimless) and diffuse (loose and weak), long-winded (tiresome) and prolix (trivial and boring), redundant (repetitious), or all of the above. See defuse, diffuse; plain English, plain language; redundancy; verbiage. Also see Garbl's Concise Writing Guide; Garbl's Plain English Writing Guide; Garbl's Fat-Free Writing Links; Garbl's Plain Language Resources.
concluded See attribution.
confidant, confident Sometimes confused. You tell secrets or intimate details to a trusted friend or confidant. If you're self-assured or certain about something, you're confident.
confute See refute.
congress Capitalize U.S. Congress and Congress when writing about the U.S. Senate and House of Representatives. Lowercase when used as a synonym for convention or in second reference to a group that uses the word as part of its formal name. Lowercase congressional unless it's part of a proper name.
congressional districts See districts.
congressman, congresswoman Use only when writing about members of the U.S. House of Representatives.
connote, denote Sometimes confused. Connote, like connotation, suggests or implies a feeling or secondary meaning besides the actual meaning of a word. Denote refers to the explicit or literal meaning of a word. See denote.
connoisseur Commonly misspelled. Double the consonants inside the word; two n's and two s's.
conscience, conscious Commonly misspelled or confused. Conscience is a noun for "a person's feelings about doing something that is morally right or wrong." Conscious is an adjective meaning "awake and aware or intended and planned."
consensus Commonly misspelled. Its first letter is the only c. Means "general agreement or opinion of all or most of the people concerned." It does not necessarily mean unanimous agreement. Avoid using the redundant consensus of opinion and general consensus. Simply use consensus or agreement. Broad consensus is acceptable.
conservation, conservative Subtle difference in emphasis, unfortunately. Conservation is an action that protects, preserves and restores works of arts, natural things like forests and wild animals, and other resources like water, gas and electricity. Conservative describes a preference for preserving established traditions or institutions and resisting or opposing any change in them--to keep doing things the traditional way despite changes in modern society. See liberal, progressive.
consolidate Try replacing with simpler combine, merge or join.
constitution Capitalize references to the U.S. Constitution, with or without the U.S. modifier: Congress is considering an amendment to the Constitution. If you're writing about the constitution of other countries or states, capitalize constitution only if the name of the country or state comes before it: the Norwegian Constitution, the country's constitution, the Oklahoma Constitution, the state constitution. Lowercase constitution in other uses: the chapter constitution. Also lowercase constitutional unless it's part of a proper name.
construct Try replacing with simpler build or erect.
consul, council, counsel Sometimes confused nouns. A consul is "an appointed government representative for aiding citizens and businesses in a foreign country." A council is "a group of people elected to represent residents of a town, city or county" and "a group of people who make decisions for an organization." Usually used as a verb meaning "to advise," a counsel is "a lawyer or group of lawyers who give legal advice and represent clients in court.
Also, council is a singular noun that should take singular verbs; the articles a or the should usually come before council. Counsel can be either singular or plural, followed by the appropriate verb form. The articles a and the are not usually needed before counsel.
Correct uses: We received the legal opinion from counsel. Counsel has suggested we go to trial on Tuesday. The company brought this matter before the council. The council advised the representatives of its position. Incorrect uses: The company brought this matter before council. Council advised the representative of their position. We go before council at noon.
contact Preferred verb meaning get in touch with or communicate with--through email, fax, telephone and postal mail. But if you mean call, write, see or similar actions, use the specific verb.
contemptible, contemptuous Sometimes confused adjectives. Use contemptible to describe something or someone that deserves contempt, scorn or lack of respect; that's despicable, worthless or disgraceful. Use contemptuous to describe a person's feelings or expressions of contempt, scorn and disdain.
content, contents These words have a subtle difference in meaning. Use content to write about the topic or subject of a book, letter, article, advertisement, speech, commercial or other written or spoken material--or to mention a single item that something contains, if necessary: Beans have a high protein content (or better: Beans have a lot of protein). Use contents to list the ingredients or items in a recipe, room, book and so on.
contiguous to Commonly misused and pompous. Does not mean "close to" or "near" but "touching and sharing a boundary." Think about using next to or bordering instead.
continual, continuous Often misused or confused. Continual means "repeatedly, often recurring or intermittent, with breaks in between": She has to repair the car continually. Periodically or intermittently are useful, clear synonyms for continually to describe something that starts and stops. Continuous means "uninterrupted, in an unbroken stream": Sales have been growing continuously for the past five years.
continue to remain Redundant and wordy. Simplify. Use remain or continue, not both, or try be still or stay.
continued Don't abbreviate. Continued, Continued on Page X, Continued from Page X, and even To be Continued are clear, concise statements. But if you must abbreviate continued for some questionable reason, use contd., without an apostrophe. Other abbreviations for continued also are abbreviations for other words.
contractions Used occasionally, contractions can speed reading and assure accuracy. They can soften the tone of your writing by making it more personal and conversational. In most writing, consider using common contractions like aren't, can't, don't, doesn't, he'll, I'll, it's, she'll, shouldn't, that's, they'll, they're, they've, you'll, you're, wasn't and won't. Avoid excessive use of contractions with dual meanings, like I'd and he'd, because they can mean both I had and I would, he had and he would. Other awkward or uncommon contractions to avoid in writing: it'd, I've got, should've, who're, would've and you'd. See could of, may of, might of, must of, should of, would of; Myths and Superstitions of Writing.
contradict See refute below.
contribute, contribute to Overstated. Simplify. Try give for contribute and add to for contribute to.
control, controlled, controlling
controversial See noncontroversial.
conversate Not a word. Replace with converse. Better still, simplify. Replace with speak, talk or even chat.
convince, persuade Often confused. Convince involves thought, trying to affect a person's point of view. Persuade involves action, trying to get a person to do something. Convince usually goes with of or that: He convinced his boss of his value to the company. She convinced her colleague that she was right. Persuade usually goes with to: The students persuaded their teacher to extend the deadline.
cooperate Think about replacing with simpler help.
copy edit, copy editing, copy editor Two words each.
copyright Sometimes misspelled as copywrite. Use the verb and noun to describe a legal right to produce, publish and sell a book, play, song, photograph, print and so on. It's not only about writing. The adjective is copyrighted.
cord See chord, cord.
corp., corporation See company names.
Corps of Engineers On first reference, use U.S. Army Corps of Engineers. Corps of Engineers is acceptable on later references.
corroborate See collaborate, corroborate.
cost-effective Jargon, cliche. Think about substituting with economical or efficient.
could See can, could.
could (not) care less If you care somewhat about something, drop the not. But if you don't care at all, keep it.
could of, may of, might of, must of, should of, would of Frequent misspellings of could have or could've, may have, might have or might've, must have, should have or should've, and would have or would've. Also, avoid using those awkward contractions in writing. See contractions.
council, counsel See consul, council, counsel.
council districts See districts.
councilmember Use the non-gender word councilmember instead of councilman or councilwoman. Capitalize only when used as a formal title before a person's name: Ellensburg City Councilmember Steven Fujita attended the meeting. Lowercase when it stands alone: The councilmember spoke at the meeting. See capitalization.
county Capitalize when part of a proper name: Clark County. Capitalize the full name of county governmental units: Clark County Personnel Department.
Always lowercase county when standing alone as a noun or used as an adjective: Population is increasing in the county. The county budget is scheduled for adoption. Lowercase plural combinations: Benton and Franklin counties. See capitalization, districts, governmental bodies.
Capitalize as part of a formal title before a name: County Executive Mary Gustafson. Lowercase when it is not part of the formal title: county Utilities Director Arnold Beck.
county council Capitalize the full name on all references: Benton County Council. Also capitalize County Council if the reference to a particular county is clear. Lowercase council when used alone: The council will meet next Thursday. Capitalize chair when used as a formal title before the name of a person in a council or committee position: Benton County Council Chair Isaac Washington. Capitalize councilmember when used as a formal title before a person's name: Benton County Councilmember Joyce Klein. Lowercase chair and councilmember when they stand alone or after a name: Kathleen Williams, a councilmember, said .... See chairman, chairperson, chairwoman; county above.
countywide One word.
couple of Follow the noun couple with the preposition of in most writing: He left a couple of style manuals in the lounge. Dropping the of and using couple as an adjective is still considered casual and slang.
course of See in the course of.
course names and numbers Capitalize the subject when used with a numeral: Geometry 2, U.S. History 101. Lowercase subjects that aren't proper names when used without a numeral: algebra, geography, Spanish.
court decisions Use numerals and a hyphen: The Supreme Court ruled 3-6, a 3-6 decision. But use the word to in direct quotations: "The court ruled 3 to 6."
court names Capitalize the full proper names of courts. Also capitalize the name if the county name, city name, state name or U.S. is dropped: Clark County Superior Court, Superior Court; Cannon Beach Municipal Court, Municipal Court; state Supreme Court, Supreme Court, U.S. Court of Appeals for the 9th Circuit, Court of Appeals. Lowercase court when standing alone. See judge.
CPR Acceptable (capitalized) in all references to cardiopulmonary resuscitation.
create, creative, creativity, creation The verb, adjective and nouns for a powerful human behavior, trait, process and act. To create is "to cause something or someone to exist" and "to produce or invent something using imagination and artistic skill." Creative describes someone who "produces or uses new and effective ideas" and "is good at using the imagination." Creativity is "the ability to use the imagination to develop, produce or use new and original ideas and things." Creation is "the act or process of inventing, producing or making something" and "something that has been invented or produced using the imagination, such as a work of art or piece of clothing." See Garbl's Creativity Resources Online.
credible, credulous Sometimes confused adjectives. Credible means "believable because evidence and logic support it." Credulous means "tending to believe too readily, gullible."
crisis, crises Sometimes misspelled, misused and overused. Crisis is singular and takes singular verbs. Crises (not crisises) is plural and takes plural verbs. A crisis is "a significant coming together of events -- a turning point -- in which the impending outcome will make a decisive or abrupt change." Avoid referring to -- and responding to -- every difficult situation as a crisis, be it an identity crisis, midlife crisis, environmental crisis, financial crisis, economic crisis or the suppposed "bankrupty" of the successful 70-year-old U.S. Social Security system nearly 40 years from now.
criteria, criterion Often confused. As the plural form of criterion, criteria is a plural noun that takes plural verbs and pronouns: The criteria are listed on the board; we will use them to test the product. Don't use the criteria is. Criterion is a singular noun that takes singular verbs and pronouns: One criterion is ease of maintenance; it is first priority for mechanics.
criticize Commonly misspelled.
criticize, critique Criticize and its various forms are becoming more negative in meaning, suggesting disapproval. Consider using critique as a neutral verb for judging both the good and bad qualities of something or someone. Its other tenses are critiqued and critiqueing, not critiqed and critiqing.
crosstown One word.
crowd See collective nouns.
cul-de-sac Always hyphenate and lowercase. Cul-de-sacs is preferred plural form.
currant, current Sometimes confused nouns. A currant is "a small round red or black berry." A current is "a flow of water or air in one direction" and "a flow of electricity through a wire."
customary Think about replacing with simpler usual if no meaning is lost.
cut and cover Hyphenate when used as an adjectival phrase: Using the cut-and-cover method was less expensive than tunneling.
cut back (v.), cutback (n. and adj.) He cut back spending. The cutback will require increased efficiency. Also, think about simplifying the verb form by dropping back.
cut off (v.), cutoff (n. and adj.) The other car cut off the truck. The cutoff date for permits is the last Friday of the month.
cutting edge, on the Cliche. Think about replacing with advanced, innovative, new, original or unconventional.
cynic, skeptic A cynic is a disbeliever. A skeptic is a doubter. Skeptics may be good journalists; cynics never are.
Maintained by Gary B. Larson of Seattle, Washington, firstname.lastname@example.org.
Updated April 14, 2013. |
Studying the Bible, obviously
Beginning to Study the Old Testament
- Peter Williams is Warden of Tyndale House, Cambridge; he was previously senior lecturer in New Testament at Aberdeen University. View all resources by Peter Williams
This article is taken from Encountering God's Word: Beginning Biblical Studies, edited by Philip Duce and Daniel Strange, and is reproduced here with kind permission of IVP/Apollos.
Preparing to Study
The Old Testament (OT) is a sizeable body of writings and, according to people of quite diverse persuasions, it was written over a very long time. The extended period of writing is itself complex, which makes study difficult. The matter is complicated further by the fact that the OT is written in a couple of obscure languages (Ancient Hebrew and a bit of Imperial Aramaic). To cap it all, the account of the origins of the OT given by many contemporary scholars is quite different from the account the OT seems to give of its own origins. So it is hardly surprising that some students beginning OT study feel a bit lost.
Yet, if the challenges of OT study can be overcome, it is also an area that yields rewards far greater than the effort put in. A good grounding in OT not only has the reward of studying the OT itself, but also pays dividends in New Testament (NT) study and in many areas of Christian theology. Moreover, the OT not only helps these disciplines but is itself a foundation for them. All of Jesus’ teaching assumes a familiarity with the OT. No one can really develop a serious understanding of the New Testament without the backdrop of the Old, nor can any Biblical or Systematic Theology afford to ignore the material of the OT. So while students choosing modules for a Religious Studies course might well be advised to delay OT study until they are ready, the delay should not be indefinite. The study of the OT is imperative. The question then is not whether to study the OT but when and how.
Answering the when question first, it is important to remember that a confused mind presented with a complicated issue is unlikely to end up clearer than it was at the start. Someone therefore approaching in-depth study of the OT will greatly benefit from prior familiarity with the OT subject matter. To put it briefly: it is helpful to have read the OT before beginning an OT course. This may seem like a high demand, but given that to read the whole OT out loud in English takes less than fifty hours, the demand is not really so high. Suggestions for reading the Bible through first time are given in an appendix.
Yet it might seem that to become familiar with the OT before studying it formally defeats the object of study. After all, the main motivation a student has in taking an OT course is usually precisely to become familiar with the OT. But at a practical level this often does not work. Most lecturers would love their students to immerse themselves in the primary text of the OT, but also need to test that students are working. They therefore set essays, and essays require arguments which in turn require the reading of secondary literature. In order to make sure that those who have read the primary text before the course do not get off too lightly lecturers will normally set for reading a quantity of secondary literature that will fill all the available time, which is generally not hard in a busy student’s life! There is similar pressure from the student’s side to read secondary rather than primary literature. There is no doubt that fifteen well-spent minutes cribbing from secondary literature will in general help a student ask smarter sounding questions than fifteen minutes spent on the primary sources. The problem is that those who pursue such a short-sighted approach over a long period are liable gradually to lose touch with the text they are supposed to be studying. Most contemporary scholarly writing will only be of minor historical interest within a generation, whereas the primary text will live on. So it is vital to read the text first. Anyone who has already started an OT course without having read the OT through should endeavour to read the whole at the first available opportunity (such as a holiday week).
This necessity of familiarity with the biblical text prior to engaging in academic study of it is in line with how we learn best. Essay topics will inevitably invite students to assess critically their own assumptions and views, as well as those of others. Students with good reasons for believing what they believe should not be anxious about such probing; it is rather an opportunity to develop mentally. However, in order to judge properly the views of authors on the bibliography it is necessary not only to evaluate the strength of the arguments they use, but also to know whether there are additional relevant contrary arguments which they do not consider. How can you critically assess any writer if you have a blank mind to start with?
So the question of when to study the OT is answered simply: study the OT when you have a firm outline in your mind of its contents and a basic familiarity with its text.
The next question is how to study, including what order to study things in. I would like here to advise strongly that someone beginning OT study should consider prioritising study of Hebrew. Not everyone finds languages equally easy or attractive, and not everyone is given the opportunity to learn Hebrew. However, this is one of the first questions you need to ask: are you going to learn Hebrew and study the OT in its original language, or are you going to study it all in translation? Here there is something analogous to the contrast above between studying primary and secondary literature. Progress with any language will seem incredibly slow at first, and much quicker comprehension of what is going on will be achieved if the same time is spent reading the English translation of the Bible. You need to weigh the effort of learning Hebrew against the length of a course in OT. If the course is three years or more then you should learn Hebrew.
The advantages of tackling original languages in Bible study are several, but only for those who persevere for the reward:
- In the long run you can save time reading secondary literature. It is still necessary and profitable, for instance, to read commentaries on books of the Bible, but you no longer need to rely on these as the major way of finding out the meaning of words.
- You can critically evaluate commentaries. A surprising number of commentaries are written by people who do not have a good grasp of Hebrew. Familiarity with the original languages will help you recognise which commentators are more worth reading and which less.
- Many parts of the Bible use specific items of vocabulary as linking words, in word play, or in other similar devices. You will be able to recognise these easily if you know the language.
- All academic disciplines are subject to fashions to one degree or another. It is possible to become immersed in the theories that prevail while one is a student only to find that they have become outdated within a generation. Language skills are less affected by fashion.
If you do decide to take up the challenge of Hebrew, you should take it up sooner rather than later. You should begin learning, and then, by seeking to use the Hebrew text at every opportunity, all subsequent study can revise your Hebrew. If you do not, then you need to get used to consulting a variety of translations of the Bible, and to learn to suspect suggestions made in commentaries about the original language that cannot be verified in at least one other source.
New textbooks for teaching Hebrew are being produced all the time and there are also less traditional helps. Each has advantages, though there is no way to avoid the fact that much time has to be invested in order to learn a language. Many recent grammars do not contain exercises of translating English into Hebrew. This is fine for a low level grasp of a language, but ultimately if you want to have a firm grasp of the language it will be necessary to undertake the rather artificial exercise of writing Biblical Hebrew.
A Christian Approach to the Old Testament
It is impossible to study without presuppositions, and to admit to having presuppositions is not to admit to having an incurable disease of brain-bias. Since no one is without presuppositions, the question is not whether one has presuppositions or not, but whether one has good presuppositions or not, and whether they are held in the right way. Presuppositions are belief systems that can direct observation and affect the way data are understood. Belief systems, whether ‘religious’ or not, can function at more than one level. An overarching system could be that there is a God who acts in the world. A lower-level system could be a particular literary or linguistic view. Individuals can sometimes share higher-level beliefs without sharing lower-level ones, or share lower-level beliefs without sharing higher-level ones.
This is an explicitly Christian introduction to OT study, and its approach is moulded by the overarching beliefs that God is truthful and that the Bible is speech from God. These statements need both elaboration and defence and space here permits neither. Nevertheless, I hope that the rest of this chapter at least provides a little of both.
To treat the OT as God’s own speech makes a difference right from the moment you open a Bible. It is not only that we probe it, but also that its contents probe and challenge us deeply. When we read that ‘The fear of the Lord [YHWH] is the beginning of knowledge’ we are faced with a choice. Either it is true, in which case the way to be knowing, even about the statement itself, is to adopt a reverent attitude before YHWH, the God revealed in the OT, or we can reject the statement without having entered into the state exhorted by the verse. There is no possible position of abstract consideration of the truth of this statement: you either adopt the attitude it encourages and fear YHWH, or you reject its claims in both thought and experience. While the YHWH-fearing cannot claim objectivity, neither can the person who has rejected the statement without ever following the exhortation. It can be argued that religious writings outside the Christian faith also exhort followers to engage with them by a similar trick of insisting that personal experience is the only test. But the point here is simple: according to claims within the OT it is not possible to consider the Lord in a dispassionate way. Involvement with the God found in the OT means more than mental assent, but life, desires, and everything. So a Christian student of the OT needs to be involved beyond the mental level. Prayer and obedience need to integrate with study. A Christian needs to adopt the habit of praying in relation to the Bible, both before reading it and in response to reading it. And what is prayed through needs to be worked out in life too. After all, according to Paul one of the primary purposes of the OT is to teach people to live righteous lives.
But there is surely a problem here. If the student is involved in this way in study will this not lead inevitably to being uncritical, and to religious bias clouding judgement? Not necessarily. Arguably those who are explicitly aware of their biases will be more able to prevent them from clouding their judgement than those who deny that they have such biases. Moreover, as one studies the OT one is also often studying the supposed basis for beliefs. Students may find with time and study that certain starting assumptions they have are not supported in the way they thought, and this can lead to modifications of belief. The key thing here is that change is made after due reflection and after one has sought at appropriate length for explanations within the framework within which one is working. A leaking roof in an individual’s world-view does not necessitate abandoning the building. Christian students encountering faith problems should be familiar enough with what has been written by those sympathetic to their view to know whether there are ways of plugging the leaks.
Most university and college courses aim to teach students to think critically by setting a number of written assignments in which a student has to evaluate contrary arguments. In order to maximise the intellectual exercise assignments are generally set in areas of controversy. Over the course of a year’s study a student may look at literally dozens of such controversial issues, often with considerable time constraints. Courses are not, however, designed so that students reach closure on an issue. That is, it is not expected that students will have reached the definitive answer by the time they complete their essays. The exercise is to start people thinking about issues. This is important in terms of how students regard conclusions they come to at the end of their essays. Essay conclusions should not be viewed as definitive results, and students need to be ready to revise their views with subsequent study. It also means that a student needs to be ready for more questions to be opened than are closed. It generally takes much less time to ask a question than to give a coherent answer. This may mean that at the end of a year there are a lot of open questions, and in OT studies some of these will inevitably touch on matters of faith. It needs to be recognised up front that the number of such questions that a student has failed to resolve bears little relationship to whether the questions are soluble or not but much relationship to the general way that people are taught to think in the academic system. Thus it is entirely plausible that someone could end up with much greater mental confusion after a course than the subject really needed to involve. This is not anyone’s fault, but simply a product of a learning system, which nevertheless has its advantages. Returning to our house analogy: everyone wants to live in a habitable house; if there are too many unresolved issues in someone’s mind then it may be like a house with a large number of leaks. There is no reason to believe that any of the leaks cannot be plugged, and there is absolutely no reason to doubt the soundness of the basic structure of the building, but the leaks may well make someone want to move out. To speak plainly, unanswered questions in large quantities may make people want to change their beliefs. However, unanswered questions in large quantities can result from two quite different sources: (1) from the inadequacy of the student’s belief system; (2) from a problem-centred method of learning. It is important to distinguish these two. The only way that this can be done is by going out of your way to look for answers to the problems that arise, and to do so persistently. It is only after trying this that one should consider any large shift in belief.
Related to this, it is important to judge correctly how certain the knowledge acquired during a course actually is. What to an undergraduate writing an essay on the basis of three days’ research seems certain is often not at all certain in real terms. It is frequently a conclusion drawn from a limited understanding based on a selective reading of the items that the student has selected to read from what a lecturer has selected from his or her already selective knowledge of writings on the topic. The writings that exist on any topic are a small subset of what could have been written on that topic, and their makeup is affected by human circumstance and market forces, and, above all, by the limits of what has yet been discovered. Humans need therefore to admit the provisional nature of their knowledge. This is difficult in an academic environment where essays are generally a means of showing off knowledge to a tutor, and where admitting ignorance is hard.
If undertaking a broadly secular course the Christian student should read a wide variety of literature from diverse standpoints. At a university or college, students of the OT will inevitably be challenged to read books by those who do not share their view of the OT. To interact with this literature is an important part of the intellectual training the course provides. However, if a course is set from a secular viewpoint it is likely that Christian students will find that the course leaves little time for them to get to know their own heritage. Of course on some topics it may be that the approach which the student finds most helpful is not written from a Christian perspective. It may be more helpful because it provokes that student, or simply because it is a fairer treatment of a subject. However, students need to beware of reading and hearing much that has been said from non-Christian points of view and yet imbibing little of what has been written from a Christian point of view. To overcome this practical hurdle requires a student to be organised and to arrange time to read Christian approaches on the same issue. Because of limitations of time during a course it may even be beneficial if a student reads such approaches before beginning the course, but since advance warning of essay topics is often not given this is difficult.
So what is read is important. Yet it is also important to know how to read. Frequently the best book on a subject is not a Christian book. Most often there simply is no ideal book on a subject from any standpoint, but several from which the relevant information needs to be gleaned. For this reason it is necessary to develop the ability to evaluate secondary literature critically, to separate information and data that an author gives from their assumptions, and to develop a sense of what needs to be double checked.
Sometimes scholars are kind enough to state their world-view. Consider the following quotation from the well-published OT scholar Philip R. Davies in an essay entitled ‘Whose History? Whose Israel? Whose Bible? Biblical Histories, Ancient and Modern’.
The belief in a single transcendental being who can comprehend, indeed controls, all history is precisely a biblical belief: it is one of the major tenets of biblical historiography... When I claim, then, that there is no ‘objective’ history I am implying a world-view incompatible with that of the biblical writings (except perhaps Qoheleth) for whom history was defined by divine deeds…’
In the ensuing context Davies criticises other scholars who have indicated that they are atheist or agnostic for adopting the Bible’s framework that there is a definitive view of the past, and thus for having a ‘theistic’ approach to history. The conclusion of the essay that a good historian needs ‘to remain sceptical, minimalist and negative’ gives us insight into his approach and will not entice many students with a faith commitment. Unfortunately Davies’ candour is rare. This means that a student needs to be aware of both the wider agendas that can influence the work of scholars and of how a Christian approach may lead to different conclusions.
One of the main areas where the historic Christian approach differs from other approaches within contemporary academia is in seeing overarching harmony within biblical writings. If biblical writings are genuinely expressions of a coherent divine mind then they must cohere at the level of what they communicate even if ultimately the nature of the coherence is beyond mortal comprehension. Thus while there is no problem in the admission that contrary statements appear in the Bible these are seen as linguistic codes for a message which at a deeper level is coherent. Contradiction is a perfectly legitimate method of communication, one frequently used by humans, but which when found in the Bible is used to deny coherence. However, the issue of coherence in the Bible is one which must be explored beyond the level of mere contrary statements, which can easily be seen to have slightly different referents.
One of the attractions, however, of secular academic methods is that, according to them, the interpreter is not bound to see a coherence between two different texts. This ‘liberates’ the interpreter from the constraint of fitting together texts which are incoherent and is held to be fairer to the text. The appeal thus of secular critical methods is that by them the interpreter is being more faithful to the Bible than within a framework ‘constrained’ by seeking to find a deep harmony between passages in Scripture. This view is a strong element motivating people to abandon a classical view of Scripture’s coherence derived from the Bible’s statements about God’s word.
Nevertheless, this appeal needs to be examined to see whether it is all that it seems. In fact, much modern research on how language works indicates the importance of unexpressed elements in the context of a text for its correct understanding. The proper interpretation of texts therefore can require a wide context. Historical-critical scholars can, however, by their focus on the true meaning of small units of text be in danger of treating small isolated texts as self-interpreting. It is important to see clearly that fairness to a text does not necessitate viewing it in isolation. What is the natural reading of a text in isolation may in fact be a distortion of its meaning, while what is a less natural reading of a text in isolation, but one motivated by reading it in a wider context, may be the correct one.
An example of this are passages in the prophets that seem to express prophetic disapproval of the OT sacrificial system and festivals (e.g. Isaiah 1:11–14). It would be possible to treat these statements as an isolated unit by some hypothetical author who believed that animal sacrifices should not be offered to God. However, viewing the verses within a wider context allows another interpretation.
But there is a danger here. It is possible so to qualify the interpretation of a text by a wider context that one in fact drowns out the message of the individual text. An interpreter can bring a text into false harmony with another text. This must be avoided while also being wary of historical-critical approaches that use the identification of incompatibly divergent voices within Scripture as a foundation upon which to build a history of OT religion. Proponents of this method can seek to discover divergent voices within the OT since such divergences are the building blocks to write a ‘critical’ history, that is one not dependent on an uncritical acceptance of scriptural statements. The way to avoid misinterpretation seems to be to read texts as a whole and simultaneously with regard to the detail they contain. Detail can modify one’s interpretation of the whole, as can the whole modify one’s interpretation of the detail.
The interpreter of Scripture also needs to be aware of the ramifications of positions taken on issues of interpretation. While coming to a view that two biblical texts cannot express truth about a coherent reality may appear like the best solution to an interpretative problem, it makes it hard to view both texts as coming from a single divine author. A student should expect a corpus as rich and diverse as the OT to present a balance of statements on a number of subjects and should not quickly conclude that variety of statements indicates a variety of theologies in the OT, if by ‘a variety of theologies’ is meant incompatible rather than merely complementary theologies.
It might be urged that to encourage students to be patient in recognising unity within Scripture will lead to an uncritical approach. There are, however, many parallels for this within other academic disciplines. For example, while a popular perception of science is that it proceeds from evidence to understanding to belief, this is very often not the case. Invariably a student of quantum mechanics believes it before understanding it, often even believing it although never expecting to understand it. The order is as in Anselm’s saying credo ut intellegam ‘I believe in order that I may understand’. For Anselm, as for Augustine before him, belief preceded understanding. This is a frequent human order, because without the belief that something makes sense it is often difficult to summon the patience to see how it makes sense. If students believed quantum mechanics to be incoherent and to have no correspondence to the universe they would certainly not have the patience to study it as intensely as if they believed that it was one of physics’ great unifying theories with wide explanatory powers. There is, of course, a danger with credo ut intellegam, since it can be used to justify continuing to believe in something unworthy of belief. However, the principle of patient positive study of Scripture to find unity is not something that should be undervalued.
One of the first issues that any student of the OT is faced with is the issue of historicity, or the question ‘Did it happen?’ Here one needs to be up-front about the importance of a historical basis for the Christian faith, and about the fact that the Church’s consensus up until the time of the Enlightenment was that the OT was true history. There were, of course, many for whom literal historicity was not the major concern. A prominent church father, Origen (ca. 185–254), was one of a number who stressed the primacy of the allegorical interpretation of Scripture. But this should not be taken to mean that he denied the general historicity of Scripture. Nor were considerations of the OT’s historicity without qualification. Martin Luther, for instance, accepted the book of Job as history but held that Job’s speeches in the book were not records of his actual words, but a creative expression by the author of Job’s thoughts. Luther thus maintained the historical nature of biblical books, but also that truth can be presented in a variety of literary genres.
It Is Literature Too
‘Genre’ has, in fact, been a much-used word recently as a number of scholars have stressed the literary qualities of the OT. Various contemporary approaches maintain that many parts of the Bible contain literary structures showing careful composition. The relevance of this to the issue of history in the OT has been to shift the focus of investigation of certain texts from their historicity to their literary nature. Authors with literary approaches range from those who use the literary nature of a work to deny its historical nature through to those who accept the historical nature of the work but choose to focus on literary structures. Some such approaches to a text are called ‘synchronic’ because they seek to view the text without regard to temporal distinction in origin, and thus interpret the whole without considering historical development. The synchronic approaches are contrasted with the diachronic ones which look at the development of texts through time. Whatever approach ultimately is taken it must be justified on a basis other than widespread use. It is important that texts are allowed to speak for themselves and that a serious attempt is made to understand the focus and themes of a text as a whole. This said, no student can afford simply to ignore historical questions.
Recognising a literary focus may even alter one’s view of the history to which a narrative refers. Thus King Jeroboam II of Israel receives relatively little treatment in the biblical narrative, much less than King Ahab, though there can be little doubt that in terms of his political achievements Jeroboam was by far the more significant. The emphasis of the biblical text thus differs from what would be given by a modern historian. The same mismatch between thematic biblical emphasis and generally applied quantitative criteria for importance is present with regard to the Babylonian exile. According to Jeremiah 52:29 the number of captives taken by Nebuchadnezzar king of Babylon at the time of the destruction of Jerusalem in 587/586 BC was 832. Thus one of the most significant events in the whole OT, round which Jeremiah and Ezekiel focus their narratives, which is the precondition for a book like Lamentations, and which is the culmination of the books of Kings (or arguably Genesis to Kings) is the deportation of a relatively small number of people. By contrast a much more significant military event just over a century earlier is passed over in these words: 'In the fourteenth year of King Hezekiah, Sennacherib king of Assyria came up against all the fortified cities of Judah and seized them.' If Sennacherib’s own account of this expedition is to believed, he took 200,150 captives at this time.
The average student with a general familiarity with the OT, but who had not studied it formally, would be quite surprised to know that there is no document outside the Bible that records Nebuchadnezzar’s capture of Jerusalem and taking of the people captive in 587/586 BC. This might even cause the student to conclude that the Bible’s narrative could not be relied on historically at this point. However, this would be because the student had confused narrative or theological importance with the sort of importance that would cause a foreign historian to record the events about the fall of Jerusalem. Besides, it would be to ignore the fact that the records of the other interested party, Nebuchadnezzar, no longer survive for the relevant period, the last surviving chronicle of his reign referring to 594/593 BC.
All this means that it is very important to read the biblical narrative with the utmost care to ensure that one is not attributing significance to an event recorded, which it does not actually have. This could be to attribute a wrong international importance to a narrative or wrong archaeological significance to an event. A biblical account must be carefully examined as a narrative before one can properly assess its historical implications.
Fact or Fiction?
When the narrative of the OT is considered as a whole it is obvious that there is generally more agreement about the historical nature of records recounting events towards the end of the OT period than about those that report earlier events. Though some would assign the whole OT to the category of story with no historical basis, most scholars believe that narrative books such as Ezra and Nehemiah are substantially historical, and that books relating to earlier periods have successively less historical substance until one reaches zero some time before the opening of the Bible in Genesis. Among scholars the cut-off point varies greatly: some deny a basis for the Bible’s picture of Solomon’s splendour, more deny any historical basis for the Exodus, and only a minority accept that if there were patriarchs, they did anything like what is attributed to them in the Bible.
Thus, broadly speaking, the more remote an event is from the present the more likely it is that its historicity will be doubted, disputed, or denied by scholars. This difference between nearer and more remote events is, of course, what would be expected if the biblical authors wrote basically towards the end of the period that the Bible is about, and possessed no super-human insight into the more remote periods. In that model, writers of Bible books were of course able to describe with moderate accuracy the time in which they lived and even preserve a certain amount of information about the time shortly before them. However, for more distant events they only had tradition, which was liable to corruption, and a certain level of imagination. Sometimes it is even said that in Bible times writers did not have a sense of history as we do now. However, differences in emphasis and historiography should not be taken to indicate a fundamental discontinuity between the approach to history then and now. Denial of a sense of continuity between attitudes to history then and now has particularly arisen with a postmodern view, which as part of a wider philosophical scheme refuses to acknowledge continuity between thought of the present and of the past.
However, the question of the historicity of earlier parts of the Bible needs also to be seen from another perspective. What was more remote from biblical writers is also more remote from us, and when considering the biblical period greater uncertainties attach to modern scholarly reconstructions of earlier as opposed to later stages in that period. While the chronology of events in the fifth century BC is generally fixed to within a year, the chronology of a millennium before is open to wider dispute, with three quite different chronologies, for example, of the first dynasty of Babylon, varying by about 120 years, and related uncertainties sometimes of around 64 years in Egyptian dating. It is true that one particular chronology of Egypt, the ‘Low Chronology’, is more widely accepted, but it is important for a student to appreciate the various levels of uncertainty in the study of different periods. The gaps in our information about the history of the period around 1,000 BC when, according to the Bible, David was king are very much greater than those around the time when Nehemiah was governor of Judah in the fifth century BC.
A further complication is the issue of antisupernaturalism (or naturalism) in historical investigation. It is possible, even common, for historians studying ancient texts to discount the possibility of narrated miracles, and to assume that only natural processes have been involved in the past. If processes other than natural ones have been at work then this naturalistic method will lead to historical distortion. For instance, in Deuteronomy 29:5 Moses points out to the Israelites how during forty years of wandering in the desert their sandals had not worn out. The narrative plainly appeals to something outside of ordinary processes, and there is the implication that if one sought for the usual evidence for sandal-wear it would not be there. There is the claim, then, of event without natural trace. If a scholar, therefore, were to argue that the wanderings did not happen on the ground of lack of evidence of wear on sandals, it would be an illegitimate argument. On the other hand, if a scholar were to argue that the wanderings did not happen on the ground of lack of evidence for bodies that perished in the desert, it would be a legitimate type of argument, because there is a legitimate expectation of evidence in this case.
Of course, scholars often do not come to the view that they should not make allowance for miraculous or non-natural events on the basis of mere prejudice. Often they first come to an opinion that the genre or date of literature in the Bible is not such as we can expect to give us reliable information about non-natural processes in the past. Nevertheless, the issue of naturalistic presuppositions in scholarly discussion is one that a Christian needs to address. A student must assess independently the likelihood that any statement about ancient history, particularly in regard to the history within the Bible has been established on a basis which he or she may want to question. This is not to encourage a dismissive attitude to the work of others, but merely to encourage an awareness of the theory-laden nature of people’s observations. Nowhere do the conditions of remoteness and of naturalistic presupposition come more to the fore than in consideration of the highly disputed first eleven chapters of Genesis. The issue is too complicated to address here, but any solution reached with these chapters must be able to be consistently applied in method to other parts of the OT, and within a coherent Christian theology.
But the pillars of scholarly consensus on the subject of OT history have not only been recently examined for assigning too great a level of certainty to scholarly reconstruction and naturalistic assumptions; in recent years scholarly consensus on the OT has been questioned from a new approach that is hyper-conservative in terms of sources admitted to discussion of history. Proponents of this approach have often been labelled ‘minimalists’. Niels Peter Lemche is a leading proponent of such a view and denies that the OT can be viewed in any way as historical in the classic sense of the word. Lemche clearly sees his approach as having advantages for a Christian viewpoint:
This new trend [i.e. the trend of which Lemche is a part] seems to be liberating the Bible from the tyranny of having to be historically accurate in the most minute detail in order to remain a Bible for Christians and Jews.
Lemche goes on to make a trenchant critique of the assumptions behind much OT study. Drawing on a variety of social-anthropological discussions Lemche claims that OT scholars have too often tried to view historical Israel with presupposed concepts of ethnicity and nationality that can be shown to be invalid. Nevertheless, the ‘minimalist’ critique of consensus historiography is not merely sociological. Lemche seeks to support his view with an examination of ancient documents seeking to establish that there is virtually no evidence for regarding the biblical narrative as a reliable source of historical information.
When Enemies Are Friends
Perhaps the easiest thing for a Christian student to do when faced with Lemche’s approach is to look for an instant refutation. Plenty of responses will be found both within mainstream scholarship and from explicitly Christian scholarship, which are firmly aligned in their opposition to minimalism. The lesson here is to see that the interplay between data and paradigm is extremely complex. Those with a common general outlook may disagree over data at the same time as those without a common outlook agree over it.
In order to consider the dynamics here we will look at a specific issue in OT history, namely the relationship between the biblical figure Shishak and an Egyptian Pharaoh called Shoshenk.
The Bible tells us that in the fifth year of Rehoboam’s reign King Shishak of Egypt came up against Jerusalem and took away a large amount of treasure from both the Temple and the royal palace (1 Kings 14:25–26). Shishak’s name in the Bible, especially in a consonantal form in which it occurs, šwšq, is naturally compared with the Pharaoh Shoshenk (low Egyptian chronology 945–924 BC), who invaded Palestine and recorded his invasion on a temple gate at Karnak in Egypt. The establishment of a time link, a synchronism, between the Bible and the Egyptian record has had important consequences in scholarship. First of all, it provides the earliest agreed record of an event recorded in the Bible which is also recorded in another document. Secondly, by the addition of the reigns of the kings of Judah within the Bible, the event can be given a date. This in turn is used to fix dates in Egyptian history, which otherwise would not be able to be calculated to the precision of a year. Shishak’s/Shoshenk’s invasion is calculated to ca. 926 or 925 BC.
This striking agreement between the Bible and an Egyptian record has been called into question by some of those with more minimalist leanings. Garbini criticises the standard reconstruction for arbitrary manipulation of Egyptian chronology in order to achieve an agreement with the Bible. He argues that if Egypt were given its most natural chronology without biblical interference it would be seen that the campaign of Shoshenk did not take place during Rehoboam’s reign, but during the reign of his predecessor Solomon. The Bible has moved it in order not to take away from the glory of Solomon’s reign. A different critique of the synchronism is made by Lemche. He argues that there is a strong distinction between the Egyptian and biblical texts. Whereas the biblical text records that Shishak came up against Jerusalem (1 Kings 14:25), the Egyptian record lists a large number of towns captured during Shoshenk’s campaign, but gives no indication that he visited Jerusalem, or even came near it. Lemche concludes that the biblical account of Shishak’s campaign to Jerusalem has been made up by a later writer on the basis of a memory of a campaign by Shoshenk in Palestine generally. The detail that Rehoboam handed over gold to get rid of Shoshenk was made up on the basis of the way Hezekiah in 701 BC handed over gold to buy off the Assyrian Sennacherib.
So far there appears to be a spectrum of opinion with minimalists at one end. A critical consensus stands between them and maximalists, among whom are a number of explicitly Christian scholars. A Christian student with a natural proclivity towards maximalism may be tempted to proceed at once to a critique of minimalism, perhaps even using arguments supplied from within an academic consensus to which they do not ordinarily belong. The problem is that this simple spectrum view does not reflect the whole truth, for this match between Egyptian history and the Bible has also been challenged from a quite different angle. As noted previously (footnote 25) there has been a small even more ‘maximalist’ trend in recent writing, which has been prepared to question more radically the foundations of Egyptian chronology. It has not succeeded in making a very significant dent in the consensus, but it has been most interesting to note some striking parallels between the writings of minimalists and these strong maximalists. Thus Bimson has challenged the equation between Shoshenk and the biblical Shishak. A key point in his argument is the geographical mismatch between the description of the campaign in the Bible and that in Karnak. Thus a scholar who has a very positive view of the historicity of biblical material is found to be using arguments with considerable affinity to those of one with a very negative view of biblical historicity. This illustrates how important it is to reject a spectrum view of modern scholarly literature. Almost all scholars have a large familiarity with material related to the OT, but their analysis can also be influenced by their wider beliefs. It is important therefore to understand a scholar’s work as a whole and to attempt to separate as far as possible information given from the interpretation put on it. Those who are positive towards the historicity of a narrative should not rule out the possibility that a minimalist’s critique of a scholar arguing for a more positive view of historicity is correct. Sometimes a scholar occupying a middle ground can be more guilty of mixing methodologies than a scholar at the extreme. This said, there is a lot of strength in the central ground. Despite critiques of Egyptian chronology the present consensus can be seen to provide considerable historical confirmation of the biblical account. The Egyptologist K.A. Kitchen, for instance, observes the vast amount of gold available in the Pharaoh’s coffers at a time directly after Shoshenk’s raid. It is not far-fetched to suppose that some of this is the gold which, according to the Bible, Rehoboam inherited from his fabulously rich father and surrendered to Shishak. But the records are mute, and coincidences like this may not provide incontrovertible proof of the correctness of a scholarly construction. So where does that leave us? A range of competent scholars come to quite different views of biblical historicity. Why do they do so? It is hard to deny that the results have at least as much to do with presuppositions as with data. Presuppositions can work at the highest level of one’s view of life as a whole, but can also provide smaller paradigms within which data are observed and ordered. A student needs to be aware of the interplay of data and interpretation, to avoid an approach which too readily claims historical confirmation of the Bible, and to see the paradigm-myopia of so many claims of disproof of the historicity of the Bible. A proper critique of minimalism needs first to understand what is right within minimalism.
There are of course many problems with minimalism. To bolster its case minimalism has resorted to charges of archaeological forgery; it is constantly having to push the interpretation of data to the limits (and beyond). Moreover, its negative presuppositions are never given full justification.
The most natural reading of quite a few sources is that they provide confirmation of the historicity of some of the events and persons within the biblical narrative. The seal impressions of a number of figures in the Bible have been found: King Ahaz, King Hezekiah, Jerahmeel the king’s son (Jeremiah 36:26), Baruch son of Neriah (Jeremiah’s scribe, in whose handwriting the book of Jeremiah would have been written), and Gemariah son of Shaphan (Jeremiah 36:12). The Mesha stele and Sennacherib’s annals provide, respectively, contrasting accounts of Mesha’s engagement with Israel, the subject of 2 Kings 3, and of Sennacherib’s siege of Jerusalem in 2 Kings 18–20. When adequate account is made of the difference of perspective of these enemy kings the outline of the biblical narratives receives broad confirmation. These documents, of course, are a long way from providing confirmation of the historicity of the Bible as a whole. In the writer’s opinion both attempts to prove the historicity of the Bible by external ‘confirmation’ and attempts to disprove its historicity have failed largely due to faulty paradigms. For this reason students in seeking to adopt positions consistent with the nature of Scripture as God’s Word need to distance discussion of detail somewhat from paradigms.
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Can Prophets Predict?
How Many Authors?
Is Genocide Ever Fair?
Psalms of Hate?
Use of Background
Towards the New Testament
Appendix: Reading the Bible Systematically
It is a great privilege to read the Bible and thereby to have access to the mind of God. Here advice is given about how to read through the Bible systematically for the first time. The key thing is to read different books differently. With 929 chapters, it would be possible to read the whole OT at three chapters per day in 310 days. However, when this is attempted for daily devotions it often leads to dry patches which, unfortunately, make many a reader give up before finishing Leviticus. This is because not all parts of the Bible are equally accessible, even though all parts have a profound message for us. Thus while the narratives of Genesis may be gripping, in our culture few find the book of Jeremiah as easy to read. However, some of the hardest parts of the Bible are those which, in the long term, will be most rewarding. All that is advised here is that at first these difficult sections are not read as the prime ingredient of devotions. Broadly speaking for a first read of the Bible it is best to begin with the historical books, that is the narrative from Genesis to Esther in the order these books are found in the English Bible. These provide the historical backbone against which the prophets from Isaiah to Malachi can be understood.
Since particular sections of the OT are often found to be hard work on first reading, the best way to read them is not as part of regular devotions, but by specially setting aside a more extended time, perhaps during a vacation, to work through them. Ideally these special reading times would occur roughly in line with the place of the books in the canon. The parts of the Bible that I would specifically advise are read in such intense sessions are:
1 Chronicles 1–9, 23–27
These chapters are no less important than those that surround them. But sometimes the major lesson from them is one made over the course of many chapters, not one summarised in a few sentences. The remaining books: Psalms, Proverbs, Ecclesiastes, and Song of Songs have a timeless quality about them which means that they can be read at any stage, though the last three are probably best read after one has read of the life of Solomon in 1 Kings. Psalms, however, is just the opposite of the books listed above. Because the psalms are largely prayers they speak directly to our own experience, and most of them are readily accessible to us. This means that if we read too much of them at any time we may end up taking less in than we would at a more leisurely pace. The best thing then with Psalms is to read the 150 psalms approximately one at a time over a more extended period.
Coming to the New Testament for a moment. All of the NT books can easily be read in single sittings or spread out over a period. It is often useful to read NT alongside OT. Chapters in the Gospels and Acts are generally longer than those of other books. Reading half a chapter of the Gospels or Acts each day or a whole chapter of the other books of the NT, one can read through the NT in 377 days, or approximately a year.
In conclusion, it is relatively simple to draw up a scheme of systematic Bible reading. But it is most advisable to allow flexibility in it for the type of material being read. This will make best use of the fact that, while much of the Bible speaks directly to our experience, other parts speak to us more on the grand scale of showing how God works.
[This appendix appeared in the original RTSF booklet]
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Deuteronomy 19:15.
These presuppositions are discussed in Paul Helm and Carl Trueman, eds., The Trustworthiness of God: Perspectives on the Nature of Scripture (Grand Rapids: Eerdmans, 2002).
Proverbs 1:7.
2 Timothy 3:16–17.
In Lester L. Grabbe, ed., Can a ‘History of Israel’ Be Written? (JSOTSS 245; Sheffield: Sheffield Academic Press, 1997) 104–122. Quotation from pp. 116–17.
Davies, ‘Whose History?’, 117 n. 19. Davies even criticises Karl Marx’s view of history as ‘theistic’.
Compare, for example, ‘I repent that I have made Saul king’ (1 Samuel 15:11) with ‘The Glory of Israel does not lie or repent’ (1 Samuel 15:29); or ‘They feared the Lord’ (2 Kings 17:33; cf. 17:41) with ‘they do not fear the Lord’ (2 Kings 17:34). If these statements were treated in isolation they could easily be ascribed to different levels in the compositional process. However, why should they not come from a single author?
Witness the opening of Charles Dickens’ classic 1859 Tale of Two Cities, ‘It was the best of times, it was the worst of times...’
This is at the heart of the linguistic discipline of Pragmatics. A useful introduction to Pragmatics is Stephen C. Levinson, Pragmatics (Cambridge Textbooks in Linguistics; Cambridge: CUP, 1983).
John F.A. Sawyer, Prophecy and the Biblical Prophets (rev. edn; Oxford: OUP, 1993) 21–24.
See for instance the argument of Paul in 1 Corinthians 15:1–19.
N.P. Lemche, The Israelites in History and Tradition (London: SPCK, 1998) 2, holds rather bizarrely that ‘the historical reading of the Bible is a comparative newcomer in comparison to such venerable procedures as allegorical understanding or typological interpretation…’. He manages to maintain this by stressing that exclusively historical interpretations only developed more recently.
Theodore G. Tappert, ed., Luther’s Works, Vol. 54: Table Talk (Philadelphia: Fortress, 1967) 79–80.
Robert Alter, The Art of Biblical Narrative (London: Allen & Unwin, 1981); Robert Alter and Frank Kermode, eds., The Literary Guide to the Bible (London: Collins, 1987); Adele Berlin, Poetics and Interpretation of Biblical Narrative (BLS 9; Sheffield: Almond, 1983); Shimon Bar-Efrat, Narrative Art in the Bible (BLS 17; JSOTSS 70; Sheffield: Almond, 1989).
This latter category includes many of the authors in Leland Ryken and Tremper Longman III, eds., A Complete Literary Guide to the Bible (Grand Rapids, MI: Zondervan, 1993).
Quite closely related to synchronic approaches is the ‘canonical’ approach, developed especially by Brevard S. Childs. See his Introduction to the Old Testament as Scripture (London: SCM, 1979). See also the evaluation of this with qualifications by Paul R. Noble, The Canonical Approach: A Critical Reconstruction of the Hermeneutics of Brevard S. Childs (BIS 15; Leiden: Brill, 1995).
Jeroboam II is treated in the short passage 2 Kings 14:23–29 whereas Ahab is treated from 1 Kings 16:29–22:40, though some of this is taken up with the figure Elijah.
2 Kings 18:13. I am of course exaggerating the case slightly. Jeremiah 52:28–30 and 2 Kings 24–25 attest more than one deportation at the time of the Babylonian exile, involving more than 832 people.
A worthwhile though tough aid to such reading would be Meir Sternberg’s The Poetics of Biblical Narrative: Ideological Literature and the Drama of Reading (Bloomington: Indiana University Press, 1985).
This is generally the approach of Lemche in The Israelites.
Speaking of Ezra and Nehemiah, David J.A. Clines says ‘From these books we learn virtually all we know about the history of the post-exilic community.’ See his Ezra, Nehemiah, Esther (Grand Rapids, MI: Eerdmans, 1984) 14.
Cyrus H. Gordon and Gary A. Rendsburg, The Bible and the Ancient Near East (4th edn; New York: W.W. Norton & Co., 1997) represent a more maximalist trend towards viewing the Bible as historical while working within a secular framework.
Thus Lemche, The Israelites, 2, quoted above.
Paul Åström, ed., High, Middle or Low? Acts of an International Colloquium on Absolute Chronology Held at the University of Gothenburg 20th–22nd August 1987, Parts 1–3 (Gothenburg: Paul Åströms Förlag, 1987–1989). Very much on the fringe of scholarship, a small group has proposed ultra-low chronologies that reduce parts of the Egyptian low chronology by several centuries during the Third Intermediate Period (normally around 1069–664 bc). Proponents claim that such schemes are alternative paradigms into which conventional data can be arranged so as to provide striking confirmation of biblical narratives. See Peter James, et al., Centuries of Darkness: A Challenge to the Conventional Chronology of Old World Archaeology (London: Pimlico, 1992) and the journalistic David Rohl, A Test of Time: The Bible—From Myth to History (London: Century, 1995). Several scholars reply to the former book in the review feature in the Cambridge Archaeological Journal 1:2 (1991) 227–53, to which the authors respond on http://www.centuries.co.uk.
One scholar who proposed a host of practical objections to the Exodus narrative was the nineteenth century figure John William Colenso, whose objections were particularly shocking at that time because he was Bishop of Natal. See his The Pentateuch and Book of Joshua Critically Examined, Volume 1 (London: Longman, Green, Longman, Roberts, & Green, 1862). Not all his arguments are equally compelling, and several are predicated on an overly pedantic reading of the text. An example of a rather obtuse objection is that he rejects the historicity of the numbers of Levites given in the Pentateuch since they are counted in the first census in the book of Numbers as 22,000 males, and in the second 38 years later as 23,000 males (p. 110). His objection, based on population growth statistics in England from 1851 to 1861, is that they should have grown more (as he estimates starting from 22,000 to 48,471 thirty-eight years later). However, many of his objections are more substantial.
After a long absence the possibility of constructing theories of origins outside of naturalistic assumptions is now firmly back on the agenda; see William A. Dembski, The Design Inference (Cambridge: CUP, 1998) and Dembski, Intelligent Design (Downers Grove, IL: IVP, 1999).
For an outline of the interaction between science and faith see Philip Duce, Reading the Mind of God (Leicester: Apollos, 1998). An approach which seeks to address some of the scientific issues but without the usual positivistic trappings of a Christian apologist is Leonard Brand, Faith, Reason & Earth History: A Paradigm of Earth and Biological Origins by Intelligent Design (Berrien Springs, MI: Andrews University, 1997).
Scholars associated with this group include Niels Peter Lemche and Thomas L. Thompson, both of Copenhagen, and Philip R. Davies of Sheffield. Lemche and Thompson prefer the term ‘maximalist’ for themselves, but this is in part due to a rhetorical technique used by the Copenhagen scholars whereby they subvert accepted meanings of terms. See BAR 23 no. 4 (July/August 1997) 28.
Lemche, The Israelites, 1.
Lemche, The Israelites, 1. T.L. Thompson’s critique of much of the academic consensus in historiography is likewise framed in rather religious terms. See his The Historicity of the Patriarchal Narratives: The Quest for the Historical Abraham (Beihefte zur Zeitschrift für die alttestamentliche Wissenschaft 133; Berlin: De Gruyter, 1974) 326–30.
Lemche, The Israelites, 35–64.
In the Bible the spelling of Shishak is normally q#y# (šyšq) but in 1 Kings 14:25 it is q#w# (šwšq), which is very close to the Egyptian consonantal spelling of Shoshenk, ššnq.
Of course, it is necessary to count the reigns of the kings of Judah back from a fixed date. For a number of dates in OT history after 853 bc Assyrian documents are compared with biblical ones to provide chronological anchors.
Giovanni Garbini, History and Ideology in Ancient Israel (tr. John Bowden; London: SCM, 1988) 29–30. There are important differences between Garbini and some other minimalists. See James Barr, History and Ideology in the Old Testament: Biblical Studies at the End of a Millennium (Oxford: OUP) 91–92.
Lemche, The Israelites, 55–57.
And we might add that according to 2 Chronicles 12:4 he captured all the fortified cities in Judah. For these see 2 Chronicles 11:5–12.
Lemche, The Israelites, 57, 187.
John J. Bimson, ‘Shoshenk and Shishak: A Case of Mistaken Identity?’, Journal of the Ancient Chronology Forum 6 (1992/93) 19–32. However, in the New Bible Commentary: 21st Century Edition, ed. D.A. Carson, R.T. France, J.A. Motyer, G.J. Wenham (Leicester: IVP, 1994) 355, Bimson identifies Shoshenk with Shishak. This may not be a revision of his earlier opinion, but rather a curtailment of discussion for a popular audience. Bimson became known for his book Redating the Exodus and Conquest (JSOTSS 5; Sheffield: JSOT, 1978), which argued that the Exodus was essentially historical.
This sort of cross-over is not uncommon. The minimalist T.L. Thompson finds himself approving of Bimson’s critique of W.F. Albright’s widely accepted chronology. See Thompson, Early History of the Israelite People: From the Written and Archaeological Sources (Leiden: Brill, 1992) 24–25.
Kitchen, ‘Where Did Solomon’s Gold Go?’ BAR 15 no. 3 (May/June 1989) 30.
Lemche’s readiness even to raise the possibility that a properly excavated inscription mentioning the ‘House of David’, or another inscription mentioning kings of the Philistine city Ekron, are fakes shows a lack of judgement. He retracted the latter opinion in a footnote. However, before examination to have raised the possibility of forgery, with its implied slur on the behaviour of others, is one of the more disingenuous traits of Lemche’s behaviour. See BAR 23 no. 4 (July/August 1997) 36–38 and Lemche, The Israelites, 182 n. 38.
This was certainly done by John Rogerson and Philip R. Davies in denying that the inscription commemorating the completion of the Siloam tunnel was from before the Babylonian exile. See Rogerson and Davies, ‘Was the Siloam Tunnel Built by Hezekiah?’, Biblical Archaeologist 59 (1996) 138–49, with a strong rebuttal in the multi-authored review feature ‘Defusing Pseudo-Scholarship: The Siloam Inscription Ain’t Hasmonean’ BAR 23 no. 2 (March/April 1997) 41–50, 68.
Take Lemche’s statement of the biblical writers that ‘Everything narrated by them may in principle be historical, but the biblical text cannot in advance be accepted as a historical source or documentation; it has in every single case to prove its status as a historical source.’ The Israelites, 29.
See for instance, Hershel Shanks, ‘Fingerprint of Jeremiah’s Scribe’, BAR 22 no. 2 (March/April 1996) 36–38.
The Mesha stele, also known as the ‘Moabite Stone’, may provide some of the earliest references outside the Bible to the dynasty of David. See André Lemaire, ‘“House of David” Restored in Moabite Inscription’ BAR 20 no. 3 (May/June 1994) 30–37. |
The race-deniers, who say there is no such thing as “race,” have a difficult time explaining why, when genetic differences of native populations across the world are mapped, the result is almost exactly the same as a map of the races. (Fig. 7-4). Thus, there is little doubt that genes differ among different populations.
All of the traits discussed in the previous chapters are caused, at least in part, by genes and, to that extent, “biology is destiny.” (Sigmund Freud). Only recently has genetics advanced to where some of the genes responsible for those traits have been identified, and only still more recently have racial differences in some of those identified genes been published. Although all humans have the same genes, the percentage of each population that has any given allele of a gene can vary from 0 (no one in the population has that allele) to 100% (everyone in the population has that allele, i.e., it is “fixed”).
It would be enlightening to present a table giving the world wide frequency of every important human allele that differs significantly between different populations, but that information is not yet available. Here are some genes for brain size and intelligence (Weiss, 1992;Plomin, 2004), behavior, skin, hair, and eyes, and diseases that are either already known to differ between populations or are very likely to differ.
The Brain and Intelligence
NBPF15 (“neuroblastoma breakpoint family, member 15,” aka MGC8902), Chromosome 1. This gene encodes multiple copies of the protein DUF1220, which is expressed in brain regions associated with higher cognitive function. Moreover, sequences of the gene are specific to different primates and, as the species become closer to humans, the number of duplicate copies increases to 212. (Popesco, 2006). Individual and racial differences in the number of copies have not yet been published.
DAB1 (“disabled-1”), Chromosome 1. This gene is involved in organizing the layers of cells in the cerebral cortex, the site of higher cognitive functions. A version of the gene has become universal in the Chinese, but not in other populations. (Williamson, 2007).
ASPM (“abnormal spindle homolog, microcephaly associated”), Chromosome 1. Its alleles affect the size of the brain; defects in theASPM gene lead to small brains and low IQ. (Evans, 2004). A new ASPM allele arose about 5800 ya in Eurasia and that allele has been suspected of increasing intelligence in Eurasia; it is common in Eurasians but absent in Africans and chimpanzees. People who speak tonal languages (e.g., Chinese) are more likely to carry two newer alleles of ASPM and MCPH1 than people in non-tonal regions. (Dediu, 2007;Mekel-Bobrov, 2005).
SSADH (“NAD(+)-dependent succinic semialdehyde dehydrogenase”), Chromosome 6. The C form increases intelligence and lifespan; the T form is 20% less efficient. (Plomin, 2004; Binghom, J., “Clever people could live 15 years longer,” Telegraph (UK), Aug. 23, 2008).
MCPH1 (“microcephaly, primary autosomal recessive 1”), Chromosome 8. The alleles of this gene, commonly called “microcephalin,” at least partly determine brain size and/or organization. (Wang, 2004). A new allele of this gene that increases intelligence arose about 37,000 ya (the confidence limit is very wide — 60,000 – 14,000 BP; Evans, 2005). This allele is common in Eurasians but rare in Africans and absent in chimpanzees.
Both the newly-discovered ASPM and microcephalin alleles were strongly selected for and spread rapidly through the Eurasian populations. These genes have been associated chronologically with two of the most revolutionary changes in human affairs – an explosion of hand-crafts in the Upper Paleolithic era (40,000 ya), and the development of sophisticated cities and the beginning of major trade routes. 1However, so far a correlation between IQ and the presence of these alleles has not been found. (Woods, 2006; Rushton, 2007a).
DCDC2 (“double cortin domain containing 2”), Chromosome 6. This gene affects the formation of brain circuits that make it possible to read. (Weiss, 2005). One allele can result in dyslexia. 2
NQO2 (“Homo sapiens quinone oxidoreductase2”), Chromosome 6. This gene clearly has effects on brain activity and might affect IQ, but that information and its population distribution are not yet published. 3
IGF2R (“insulin-like growth factor 2 receptor”), Chromosome 6. This was the first gene discovered for intelligence; possession of one of the alleles of this gene increases IQ by about 4 points. (Chorney, 1998).
DTNBP1 (“dystrobrevin binding protein 1”), Chromosome 6. It is associated with schizophrenia and has recently been linked to intelligence. (Burdick, 2006).
CHRM2 (“cholinergic receptor, muscarinic 2”), Chromosome 7, activates signaling pathways in the brain; some alleles can increase IQ 15 to 20 points. (Dick, 2007; Gosso, 2006).
FoxP2 (“forkhead box P2”), Chromosome 7. This gene affects language skills, including grammar, as well as IQ. Although many animals also have the gene, humans acquired an allele within the last 200,000 yrs that was strongly selected because the superior communications and creativity it made possible were a major advantage.
EMX2 (“Empty spiracles-like protein”), Chromosome 10, codes for the development of the cortex into specialized areas. Mismatched areas lower performance. (Leingärtner, 2007).
FADS2 (“fatty acid desaturase 2”), Chromosome 11, is involved in processing omega 3 fatty acids to produce nutrients for the brain. An allele of this gene raises the IQ of children by about 6 to 10 IQ points if they are breast-fed. (Caspi, 2007).
DARPP-32 (“dopamine- and cyclic AMP-regulated phosphoprotein”), Chromosome 17. One allele of this gene optimizes the brain’s thinking circuitry, but increases the risk of schizophrenia. (Meyer-Lindenberg, 2007).
MAPT (“microtubule-associated protein tau”), Chromosome 17. Mutations in this gene can cause neurodegenerative disorders. The H2 haplotype of this gene may have come from the Neanderthals. (Hardy, 2005). Also, physicist and mathematician Roger Penrose proposed that consciousness is a quantum effect that arises in these microtubules. (Shadows of the Mind, 1996).
PDYN (“prodynorphin”), Chromosome 20. It codes for a precursor molecule for neuropeptides, which affects perception, behavior, and memory. (Balter, 2005).
HAR1 RNA (“human accelerated region 1”), Chromosome 20. This gene codes for an RNA protein that develops neurons in the neocortex of the brain. This gene is different in the brains of humans and chimpanzees and is rapidly evolving in humans. (Pollard, 2006). Also see HAR1F, which is active in special cells that appear early in embryonic development and help form the human cerebral cortex; HAR1produces RNA that does not produce protein. (Smith, K., 2006; Pollard, 2006).
EST00083 (“expressed sequence tag”) is an mtDNA polymorphism found more often in high IQ groups. It is particularly common in Europe (less so in Asia), where it is associated with a lineage that dates back 35,000 yrs. (Thomas, 1998).
PER2 (period homolog 2, Dosophila), Chromosome 2, “is a key component of the mammalian circadian clock machinery.” “[A] high and significant difference in the geographic distribution of PER2 polymorphisms was observed between Africans and non-Africans.” (Cruciani, 2008)
ADH (“alcohol dehydrogenase”), Chromosome 4. Mutations in this gene cause Asians to have a more intense response to alcohol, including facial flushing. (Duranceaux, 2006).
PAX6 (“paired box gene 6”), Chromosome 11, controls development of the iris. A mutation of this gene is linked to impulsiveness and poor social skills, which is discernable by the appearance of the iris. (Larsson, 2007).
DRD4 (“dopamine receptor D4”), Chromosome 11, controls sex drive. (Zion, 2006). Some studies found that an allele is associated with novelty-seeking personality traits in two European populations (Benjamin, 1996), but other studies did not confirm this.
ACTN3 (“alpha-actinin-3”), Chromosome 11, codes for fast twitch muscle fibers. The R allele encodes a functional copy of the protein but the X allele does not produce the protein; 25% of Asian populations are deficient, 18% of Europeans, but less than 1% of the African Bantu population. (Yang, 2003).
AVPR1a (arginine vasopressin 1a receptor), Chromosome 12, influences social bonding and altruism in humans and some animals. People with a long promoter of the RS3 allele are more altruistic than persons with a short promoter. (Knafo, 2007).
ACE (“angiotensin I-converting enzyme”), Chromosome 17. It converts angiotensin I to angiotensin II, but is also involved in athletic ability. Racial differences are not yet known.
MAOA (“monoamine oxidase A”), X Chromosome. This gene codes for an enzyme which sits on mitochondrial membranes in neurons and degrades several important neurotransmitters, including several believed to be important in the regulation of aggression and impulsivity. (Moran, 2006). People with the short version of MAOA were found to be more violent and generally more antisocial than those with the long version. Also, people with low levels of the enzyme who were mistreated as children have significantly higher crime rates. (Moffitt, 2005; Meyer-Lindenberg, 2006). Different ethnic groups have different alleles. (Wikipedia, “Monoamine Oxidase”).
Skin, Hair,& Eyes
EDAR (“ectodysplasin A receptor”), Chromosome 2, controls hair thickness. East Asians have two copies of an allele that gives them thick hair. (Am. Soc. of Human Gen., Annual Meeting, Oct. 23-27, 2007).
MATP (“melanoma antigen transporter protein”), Chromosome 5, affects skin color. “The L374F mutation was present at an allele frequency as high as 0.96 in the German population, whereas it was completely absent in the Japanese population.” (Yuasa, 2004). There are at least 118 genes associated with skin pigmentation (Lao, 2007).
AIM1 (“absent in melanoma 1”), Chromosome 6, influences skin color. The 272K allele is common in Asian populations, such as Chinese (43.4%), Sinhalese (20.4%), and Tamils (12.1%), but is rare in Europeans (2.5%), Xhosans (Bushmen, 3.4%), and Ghanaians (4.1%). The 374F allele is exclusively found in Europeans (91.6%), but not in the other five populations (0%–1.9%). (Soejima, 2006).
TYR (“Tyrosinase”), Chromosome 11. This gene and the MATP gene have a predominant role in the evolution of light skin in Europeans but not in East Asians, who evolved light skin independently. (Norton, 2006).
KITLG (“KIT legand”), Chromosome 12. About 20% of the differences in pigmentation between people of African and northern European descent is due to different alleles of this gene. (Miller, 2007).
OCA2 (“oculocutaneous albinism II”), Chromosome 15. This gene can cause albinism, but the genetics are different in Caucasians and African Americans. (Lee, 1994). It also affects eye color. (Duffy, 2007).
HERC2, (“HECT domain and RCC1-like domain-containing protein 2”), Chromosome 15, can reduce the production of dark pigment (melanin) by adjacent gene OCA2, resulting in blue eyes, blond hair, and light skin; 97% of blue-eyed people have the same allele. The high frequency of the blue-eyed allele in Scandinavia implies that allele significantly increased reproductive success. (Eiberg, 2008).
SLC24A5 (“solute carrier family 24, member 5,” aka the “golden pigmentation gene”), Chromosome 15. An allele of this gene that changes a single amino acid in a protein plays a major role in giving Eurasians lighter skin than Africans. (Lamason, 2005). The European allele is not the same as the Asian allele. (Norton, 2006). This gene is also expressed in the brain. 4
MC1R (“melanocortin-1 receptor”), Chromosome 16. There are over thirty alleles for this gene. The gene helps determine hair and skin color, but not eye color. (Mueller, 2006). Africans (and tropical indigenous people in general) have an ancestral allele for this gene and only synonymous alleles (i.e., alleles that code for the same amino acids) of this gene; the alleles are ancient and code for eumelanin, which results in black skin and hair. (Harding, 2000). Europeans have alleles for blond, red, brown, and black hair.
KRT41P, aka KRTHAP1 (“keratin 41 pseudogene”), 5 Chromosome 17. This gene is present in chimpanzees, gorillas, and man, and codes for body hair. It was turned off in man about 240,000 ya. (Klein, 2002, p. 203).
EYCL1 (“eye color 1” aka “gey”), Chromosome 19, codes for green and blue eye color; EYCL2 (“bey1”), Chromosome 15, codes for brown eyes, and EYC3 (“bey2”), Chromosome 15, codes for brown and blue eyes. (Wikipedia, “Eye Color“). Five to ten genes may be involved in eye color.
ASIP (“agouti signaling protein”), Chromosome 20. The 8818G allele is associated with darker skin color in Africans and African Americans; since the allele also is found in African apes, it is “ancestral” in Africans. (Norton, 2006,).
Health & Disease
LCT (“lactase gene”), Chromosome 2, codes for lactase, an enzyme that catalyzes the digestion of lactose, milk sugar. An allele that enables adults to digest milk sugar arose in northern Europe only recently, between 5480 BC and 5000 BC. The allele was strongly selected and its possession by over 90% of northern Europeans may help explain how Indo-Europeans were able to spread so suddenly about 4000 ya. The vast majority of Asians and Africans do not have it, but the Tutsis more recently independently evolved a lactose-tolerant allele. (Burger, 2007). Since all children are lactose-tolerant and most adults are not, “lactose tolerance may be considered a form of neoteny.” (Wikipedia, “Lactose Intolerance”).
CCR5 (“chemokine (C-C motif) receptor 5”), Chromosome 3. The delta 32 deletion of this gene appeared more than 5,000 ya in southern Finland and may have provided some protection against smallpox. Today, only a small percentage of Europeans have this deletion (1%, though 10% of European Jews have it), but it protects them from the AIDS virus (Zimmer, 2001, p. 222-225), though it increases their risk of illness from flaviviruses, such as West Nile virus; it is not found in Asians or Africans. (Smith, 1997; Stephens, 1998).
PDE4 (“pyridoxine-dependent epilepsy”), Chromosome 5. An allele of this gene is involved in cardiovascular disease and lung cancer susceptibility. Blacks who smoke up to a pack a day are far more likely to develop lung cancer than whites who smoke similar amounts. Blacks may have less protection against lung cancer because they were subjected to less smoke, as fire is not needed as much in the tropics. (Garte, 2001).
CYP3A5 (“cytochrome”), Chromosome 7, acts to retain salt in the kidneys. It is common in Africans, who live in a hot climate where salt is lost through sweat and is not easily available. The CYP3A5*3 allele, which is non-functional, is far more common in Eurasians (96% for the Basques in the Pyrenees Mountains) than in Africans (6% in Nigeria). Thus, Africans who live in white civilizations retain too much salt, leading to cardiovascular problems. Another gene, AGT M235, which is also involved in salt retention, has a similar distribution. (Thompson, 2004; Roy, 2005).
CASP12 (“cysteinyl aspartate proteinase”), Chromosome 9. Having the non-functional version of this gene better prevents sepsis (infection of the blood and tissues by bacteria). The loss of function occurred 51,000 to 74,000 ya. (Wang, X., 2006). This gene HBB(“hemoglobin beta chain”) on Chromosome 11, codes for the beta strand of hemoglobin. A single copy of an allele of this gene protects against malaria, but two copies cause sickle cell anemia; 6 it is found mostly in people living in malarial regions of Africa and India.
CD4 (“cell development”), Chromosome 12. The 7R allele was probably very ancient in Neanderthals, but may be only 30,000 yrs old in Hss. It is a receptor for HIV. (Hanna, 1989).
BRCA1 (“breast cancer”), Chromosome 17. This gene has an allele that is involved in breast cancer. Of Ashkenazi Jewish women, 1 in 40 carries alleles of the BRCA1 and the BRCA2 gene that give them a 4 out of 5 chance of having breast cancer.
LTA4H (“leukotriene A4 hydrolase”), Chromosome 17. An allele of this gene increases the risk of a heart attack in African Americans by more than 250%, but only by 16% in whites and Asians. The gene boosts inflammation as a way to fight infections and is generally not found in Africans. Although 30% of whites have the allele, they have evolved other genes to counteract it, but the 6% of the African Americans, who acquired it by breeding with whites, have not. (Helgadottir, 2006).
APOH (“apolipoprotein H”), Chromosome 17. This gene is a major autoantigen for the production of antiphospholipid antibodies (APA) in autoimmune diseases. The APOH*3B allele is present only in blacks and is identical to the wild type APOH in chimpanzees. (Kamboh, 2004).
NOS2 (“nitric oxide synthase”), Chromosome 17, encodes an enzyme that produces nitric oxide. An allele possessed by Africans in malaria areas causes increased production of nitric oxide, which protects against the symptoms of the disease. Caucasians do not have that allele. (Keller, 2004).
CNDP1 (“carnosine dipeptidase 1”), Chromosome 18. A trinucleotide repeat sequence on this gene protects Caucasian Europeans, white Americans, and Arabs, but not blacks, from diabetic end-stage kidney failure. (Freedman, B.I., 2007).
APOE (“apolipoprotein E”), Chromosome 19. This gene plays a role in transporting cholesterol and is involved in Alzheimer’s disease. It is possible that some people may not have this gene at all which, if true, would raise some interesting questions. (Miller, 2006).
PDHA1 (“pyruvate dehydrogenase (lipoamide) alpha 1″), X Chromosome. The tree for this gene is estimated as 1.86 mya and the split between Africans and non-Africans as 200,000 yrs. There are no haplotypes shared between the Africans and the non-Africans and one site (544) is fixed in the non-African lineage (i.e., every non-African tested has the same allele, which suggests it is advantageous and ancient). (Harris, 1999.).
The reader may have noticed that genes that code for one trait may affect other, seemingly unrelated traits (e.g., PAX6, CCR5, andPAX6) and that some alleles (“ancestral” alleles) are found in blacks and chimpanzees, but not other races (NQ02, ASIP, APOH*3B,MC1R) or, vice versa, (ASPM, MCPH1).
Men and women differ by only a single chromosome (Y in men, X in women), yet the differences in that chromosome extensively affect their anatomy, physiology, and behavior. Figure 13-1 (Yang, 2006) shows how genes are expressed in the livers of female (top) versus male mice. Red corresponds to more gene expression, green to less. Even though one might think that the differences between males and females would be limited to reproduction-related differences on the X and Y chromosomes, this map shows that the differences have a large effect on genes that are expressed in the liver, which has little to do with reproduction. Thus, we should not be surprised if racial differences in genes affect much more in the body than the obvious differences in appearance.
At the present time, studies of racial genetic differ-ences have been mostly limited to mtDNA and coding nuclear DNA. Yet humans have more “junk” DNA than any other animal, and the functions of junk DNA are just beginning be discovered. Important racial differences can also be expected to be found in it as well, in the number of copies of genes, and in the gene regulators, the genetically-inherited “switches” that determine whether and when a gene is read. |
76 DEFINITIONS OF THE SIGN BY C.S.PEIRCE.
DATED TEXTS (or texts dated by M.Robin)
Representation is anything which is or is represented to stand for another and by which that other may be stood for by something which may stand for the representation.
Thing is that for which a representation stand prescinded from all that can serve to establish a relation with any possible relation.
Form is that respect in which a representation stands for a thing prescinded from all that can serve as the basis of a representation, therefore from its connection with the thing.
2 - 1867 - C.P. 1-554 - On a new list of categories .
[...] every comparison requires, besides the related thing, the ground, and the correlate, also a (mediating representation which) (represents the relate to be a representation of the same correlate) (which this mediating representation itself represents). Such a mediating representation may be termed an (interpretant), who says that a foreigner says the same thing which he himself says.
3 - 1868 - C.P. 5-283 - Consequences of four incapacities .
[...] Now a sign has, as such, three references : first, it is a sign to some thought which interprets it; second, it is a sign for some object to which in that thought it is equivalent, third, it is a sign, in some respect or quality, which brings it into connection with its object. Let us ask what the three correlates are to which a thought-sign refers.
4 - 1873 - MS 380 - Of logic as a study of signs .
A sign is something which stands for another thing to a mind. To it existence as such three things are requisite. On the first place, it must have characters which shall enable us to distinguish it from other objects. In the second place, it must be affected in some way by the object which it signified or at least something about it must vary as a consequence of a real causation with some variation of its object.
5 - 1873 - C.P. 7-356 - Logic. Chapter 5 .
Let us examine some of the characters of signs in general. A sign must in the first place have some qualities in itself which serve to distinguish it, a word must have a peculiar sound different from the sound of another word; but it makes no difference what the sound is, so long as it is something distinguishable. In the next place, a sign must have a real physical connection with the thing it signifies so as to be affected by that thing. A weather-cock, which is a sign of the direction of the wind, must really turn with the wind. This word in this connection is an indirect one; but unless there be some way or other which shall connect words with the things they signifie, and shall ensure their correspondance with them, they have no value as signs of those things. Whatever has these two characters is fit to become a sign. It is at least a symptom, but it is not actually a sign unless it is used as such; that is unless it is interpreted to thought and addresses itself to some mind. As thought is itself a sign we may express this by saying that the sign must be interpreted as another sign. [...]
A representation is an object which stands for another so that an experience of the former affords us a knowledge of the latter. There are three essential conditions to which every representation must conform. It must in the first place like any other object have qualities independent of its meaning. It is only through a knowledge of these that we acquire any information concerning the object it represents.[...] In the second place, representation must have a real causal connection to its object. [...] In the third place, every representation addresses itself to a mind. It is only in as far as it does it that it is a representation. The idea of the representation itself excites in the mind another idea and in order that it may do this it is necessary that some principle of association between the two ideas should already be established in that mind. [...]
7 - 1885 - 3-360 - On the algebra of logic .
A sign is in a conjoint relation tothe thing denoted and to the mind. If this triple relation is not of a degenerate species, the sign is related to its object only in consequence of a mental association, and depend upon a habit. Such signs are always abstract and general, because because habits are general rules to which the organism has become subjected. They are, for the most part, conventional or arbitrary. They include all general words, the main body of speech, and any mode of conveying a judgement. For the sake of brevity I will call them tokens.
[...] Indeed, representation necessary involves a genuine triad. For it involves a sign, or representamen, of some kind, inward or outward, mediating between an object and an interpreting thought. [...]
A sign, or representamen, is something which stands to somebody for something in some respect or capacity. It addresses somebody, that is, creates in the mind of that person an equivalent sign or perhaps a more developed sign. That sign which it creates I call the interpretant of the first sign. The sign stands for something, its object. It stands for that object, not in all respects, but in reference to a sort of idea, which I have sometimes called the ground of the representamen. [...]
[...] A very broad and important class of triadics characters [consist of] representations. A representation is that character of a thing by virtue of which, for the production of a certain mental effect, it may stand in place of another thing. The thing having this character I term a representamen, the mental effect, or thought, its interpretant, the thing for which it stands, its object.
[...] A sign is only a sign in actu by virtue of its receiving an interpretation, that is, by virtue of its determining another sign of the same object. This is as true of mental judgments as it is of external signs.[...]
Anything which determines something else (its interpretant) to refer to an object to which itself refers (its object) in the same way, the interpretant becoming in turn a sign, and so on an infinitum.
No doubt, intelligent consciousness must enter into the series. If the series of successive interpretants comes to an end, the sign is thereby rendered imperfect, at least. If, an interpretant idea having been determined in an individual consciousness it determines no outward sign, but that consciousness becomes annihilated, or otherwise loses all memory or other significant effect of the sign, it becomes absolutely undiscoverable that there ever was such an idea in that consciousness; and in that case it is difficult to see how it could have any meaning to say that that consciousness ever had the idea, since the saying so would be an interpretant of that idea.
[...] Genuine mediation is the character of a Sign. A sign is anything which is related to a Second thing, its Object, in respect to a Quality, in such a way as to bring a Third thing, its Interpretant, into relation to the same Object, and that in such a way as to bring a Fourth into relation to that Object in the same form, ad infinitum. If the series is broken off, the Sign, in so far, falls short of the perfect significant character. It is not necessary that the Interpretant should actually exist. A being in futuro will suffice.
On the definition of Logic.
Logic will here be defined as formal semiotic. A definition of a sign will be given which not more refers to human thought than does the definition of a line as the place with a particle occupies, part by part, during a lapse of time. Namely, a sign is something, A, which brings something, B, its interpretant sign determined or created by it, into the same sort of correspondence with something, C, its object, as that in which itself stand to C. It is from this definition, together with a definition of "formal", thah I deduce mathematically the principles of logic. [...]
A sign, or Representamen, is a First which stands in such a genuine triadic relation to a Second, called its Object, as to be capable of determining a Third, called its Interpretant, to assume the same triadic relation to its Object in which it stand itself to the same Object. The triadic relation is genuine, that is its three members are bound together by it in a way that does not consist in any complexus of dyadic relations. That is the reason the Interpretant, or Third, cannot stand in a mere dyadic relation to the Object, but must stand in such a relation to it as the Representamen itself does. Nor can the triadic relation in which the third stands be merely similar to that in which the First stands, for this would make the relation of the THird to the First a degenerate Secondness merely. The Third must indeed stand in such a relation, and thus be capable of determining a Third of its own; but besides that, it must have a second triadic relation in which the Representamen, or rather the relation there of to its Object, shall be its own (the Thrid's) Object, and must be capable of determining a Third to this relation. All ths must be equally be true of the Third's Third and so on endlessly; and this, and more, is involved in the familiar idea of a Sign; and the term Representamen is here used, nothing more is implied. A Sign is a Representamen with a mental Interpretant.
Possibly there may be Representamens that are not Signs. Thus, if a sunflower, in turning towards the sun, becomes by that very act fully capable, without further condition, of reproducing a sunflower which turns in precisely way toward the sun, and of doing so with the same reproductive power, the sunflower would become a Representamen of the sun. But thought is the chief, if not the only, mode of representation.
A Sign does not function as a sign unless it be understood as a sign. It is impossible, in the present state of knowledge, to say, at once fully precisely and with a satisfactory approach to certitude, what is to understand of a sign. ..., it does not seem that conciousness can be considered as essential to the understanding of a sign. But what is indispensable is that there should, actually or virtually, bring about a determination of a sign of the same object of which it is itself a sign. This interpreting sign, like every sign, only functions of a sign so for as it again is interpreted, that is, actually or virtually, determines a sign of the same object of which it is itself a sign. Thus there is a virtual endless series of signs when a sign is understood; and a sign neveer understood can hardly be said to be a sign.
17 - 1903 - C.P. 1=53B- - Lowell Lectures: Lecture III, vol. 21, 3d Draught .
Every sign stands for an object independent of itself; but it can only be a sign of that object in so far as that object is itself of the nature of a sign or thought. For the sign does not affect the object but is affected by it; so that the object must be able to convey thought, that is, must be of the nature of thought or a sign. [...]
[...] Now a sign is something, A, which denotes some fact or object, B, to some interpretant thought, C.
[...] In the first place, as to my terminology I confine the word representation to the operation of a sign or its relation to the object for the interpreter of the representation. The concrete subject that represents I call a sign or representamen. I use these two words, sign and representamen, differently. By a sign I mean anything which conveys any definite notion of an object in any way, as such conveyers of thought are familiarly known tous. Now I start with this familiar idea and make the best analysis I can of what is essential to a sign, and I define a representamen as being whatever that analysis applies to. [...]
My definition of a representamen is as follow:
A REPRESENTAMEN is a subject of a triadic relation TO a second, called its OBJECT, FOR a third, called is INTERPRETANT, this triadic relation being such that the REPRESENTAMEN determines its interpretant to stand in the same triadic relation to the same object for some interpretant.
The mode of being of a representamen is such that it is capable of repetition.[...] This repetitory character of the
representamen involves as a consequence that it is essential to a representamen that it should contribute to the determination of another representamen distinct from itself. [...] I call a representamen which is determined by another representamen, an interpretant of the latter. Every representamen is related or is capable of being related to a reacting thing, its object, and every representamen embodies, in some sense, some quality, which may be called its signification, what in the case of a common name J.S. Mill call its connotation, a particularly objectionable expression.
A Representamen is the First Correlate of a triadic relation, the Second Correlate being termed its Object, and the possible Third Correlate being termed its Interpretant, by which triadic relation the possible Interpretant is determined to be the First Correlate of the same triadic relation to the same Object, and for some possible Interpretant. A Sign is a representamen of which some interpetant is a cognition of a mind. Signs are the only representamens that have been much studied.
[...] As we know a sign, it is something which represents the real Truth, in some aspect of it, to somebody; that is, determines a knowledge of that Truth. This knowledge is itself of the nature of a sign. In its more perfect forms, it involves consciousness, or a representation in the conscious sign of itself to itself, somewhat as a map covering a country may represent itself. But knowledge is nothing, quite nothing but a counterfeit unless it would under some circumstances, determine conduct. It must have real effects. In fact any outward sign must, not merely as a thing, but as a sign produce physical effects in order to be communicated. [...]
[...] A sign is intended to correspond to a real thing, or fact, or to something relatively real; and this object of the sign may be the very sign itself, as when a map is precisely superposed upon that which it maps. [...] A sign is also intended to determine, in a mind or elsewhere, a sign of the same object; and this interpretant of the sign may be the very sign itself; but as a general rule it will be different. [...]
A sign is supposed to have an object or meaning, and also to determine an interpretant sign of the same object. It is convenient to speak as if the sign originated with an utterer and determined its interpretant in the mind of an interpreter.
[...] Conversely, every thought proper involves the idea of a triadic relation. For every thought proper involves the idea of a sign. Now a sign is a thing related to an object and determining in the interpreter an interpreting sign of the same object. It involves the relation between sign, interpreting sign, and object. There is a threefold distinction between signs, which is not in the least psuchological in its nature, but is purely logical, and is of the atmost importance in logic.
I call that which represents, a representamen. A Representation is that relation of the representamen to its object which consists in it determining a third (the interpretant representamen) to be in the same relation to that object.
[...] In its genuine form, thirdness is the triadic relation existing between a sign, its object, and the interpreting thought, itself a sign, considered as constituting the mode of being of a sign. A sign mediates between the interpretant sign and its object. Taking sign in its broadest sense, its interpretant is not necessarily a sign. [...]
A sign therefore is an object which is in relation to its object on the one hand and to an interpretant on the other, in such a way as to bring the interpretant into a relation to the object, corresponding to its own relation to the object. I might say similar to its own for a correspondence consist in a similarity; but perhaps correspondence is narrower.
[...] It is difficult to define a sign in general. It is something which is in such a relation to an object that it determines, or might determine, another sign of the same object. This is true but considered as a definition it would involve a vicious circle, since it does not say what is meant by the interpretant being a "sign" of the same object. However, this much is clear ; that a sign has essentially two correlates, its object and its possible Interpretant sign. Of these three, Sign, Object, Interpretant, the sign as being the very thing under consideration is Monadic, the object is Dyadic, and the Interpretant is Triadic. We therefore look to see, whether there be not two Objects, the object as it is in itself (the Monadic Object), and the object as the sign represents it to be (the Dyadic Object). There are also three Interpretants; namely, 1°, the Interpretant considered as an independent sign of the Object, 2°, the Interpretant as it is as a fact determined by the Sign to be, and 3° the Interpretant as it is intended by, or is represented in, the Sign to be. [...]
So then anything (generally in a mathematical sense) is a priman (not a priman element generally) and we might define a sign as follows:
A "sign" is anything, A, which,
(1) in addition to other characters of its own,
(2) stands in a dyadic relation Þ, to a purely active correlate, B,
(3) and is also in a triadic relation to B for a purely passive correlate, C, this triadic relation being such as to determine C to be in a dyadic relation, µ, to B, the relation µ corresponding in a recognized way to the relation Þ.
In the which statement the sense in which the words active and passive are used is that in a given relationship considering the various characters of all or some of the correlates with the exclusion of those only which involve all the correlates and are immediately implied in the statement of the relationship, none of those which involve only non-passive correlates will by immediately essential necessity vary with any variation of those involving only passive correlates; while no variation of characters involving only non-active elements will by immediately essential necessity involve a variation of any character involving only active elements. And it may be added that by active-passive is meant active and passive if the entire collection of correlates excluding the correlates under consideration be divided into two parts and one part and the other be alternately excluded from consideration; while purely active or passive means active or passive without being active-passive.
This definition avoids the niceties for the sake of emphasizing the principal factors of a sign. Nevertheless, some explanations may be desirable. But first for the terminology. I use "sign" in the widest sense of the definition. It is a wonderful case of an almost popular use of a very broad word in almost the exact sense of the scientific definition. [...]
I formerly preferred the word representamen. But there was no need of this horrid long word. [...]
My notion in preferring "representamen" was that it would seem more natural to apply it to representatives in legislatures, to deputies of various kinds, etc... I admit still that it aids the comprehension of the definition to compare it carefully with such cases. But they certainly depart from the definition, in that this requires that the action of the sign as such shall not affect the object represented. A legislative representative is, on the contrary, expected in his functions to improve the condition of this constituents; and any kind of attorney, even if he has no discretion, is expected to affect the condition of his principal. The truth is I went wrong from not having a formal definition all drawn up. This sort of thing is inevitable in the early stages of a strong logical study; for if a formal definition is attempted too soon, it will only shackle thought. [...]
I thought of a representamen as taking the place of the thing; but a sign is not a substitute. Ernst Mach has also fallen into that snare.
[...] A sign is plainly a species of medium of communication and medium of communication is a species of medium, and a medium is a species of third.[...]
A medium of communication is something, A, which being acted upon by something else, N, in its turn acts upon something, I, in a manner involving its determination by N, so that I shall thereby, through A and only through A, be acted upon by N. [...] A Sign, on the other hand, just in so far as it fulfill the function of a sign, and none other, perfectly conforms to the definition of a medium of communication. It is determined by the object, but in no other respect than goes to enable it to act upon the interpreting quasi mind ; and the more perfectly it fulfill its function as a sign, the less effect it has upon that quasi-mind other than that of determining it as if the object itself had acted upon it. [...]
It seems best to regard a sign as a determination of a quasi-mind; for if we regard it as an outward object, and as addressing itself to a human mind, that mind must first apprehend it as an object in itself, and only after that consider it in its significance; and the like must happen if the sign addresses itself to any quasi-mind. It must begin by forming a determination of that quasi-mind, and nothing will be lost by regarding that determination as the sign.
I use the word "Sign" in the widest sense for any medium for the communication or extension of a Form (or feature). Being medium, it is determined by something, called its Object, and determines something, called its Interpretant or Interpretand. But some distinctions have to be borne in mind in order rightly to understand what is meant by the Object and by the Interpretant. In order that a Form may be extended or communicated, it is necessary that it should have been really embodied in a Subject independently of the communication; and it is necessary that there should be another subject in which the same form is embodied only in consequence of the communication. The Form, (and the Form is the Object of the Sign), as it really determines the former Subject, is quite independent of the sign; yet we may and indeed must say that the object of a sign can be nothing but what that sign represents it to be. Therefore, in order to reconcile these apparently conflicting Truths, it is indispensible to distinguish the immediate object from the dynamical object.
The same form of distinction extends to the interpretant; but as applied to the interpretant, it is complicated by the circumstance that the sign not only determines the interpretant to represent (or to take the form of) the object, but also determines the interpretant to represent the sign. Indeed in what we may, from one point of view, regard as the principal kind of signs, there is one distinct part appropriated to representing the object, and another to representing how this very sign itself represents that object. The class of signs I refer to are the dicisigns. In "John is in love with Helen" the object signified is the pair, John and Helen. But the "is in love with" signifies the form this sign represents itself to represent John and Helen's Form to be. That this is so, is shown by the precise equivalence between any verb in the indicative and the same made the object of "I tell you". "Jesus wept" = "I tell you that Jesus wept".
First, an analysis of the essence of a sign, (stretching that word to its widest limits, as anything witch, being determined by an object, determines an interpretation to determination, through it, by the same object), leads to a proof that every sign is determined by its object, either first, by partaking in the characters of the object, when I call the sign an Icon; secondly, by being really and in its individual existence connected with the individual object, when I call the sign an Index; thirdly, by more or less approximate certainty that it will be interpreted as denoting the object, in consequence of a habit (which term I use as including a natural disposition), when I call the sign a Symbol.
[...] That thing which causes a sign as such is called the object (according to the usage of speech, the "real", but more accurately, the existent object) represented by the sign : the sign is determined to some species of correspondence with that object.[...]
For the proper significate outcome of a sign, I propose the name, the interpretant of the sign. [...]
Whether the interpretant be necessarily a triadic result is a question of words, that is, of how we limit the extension of the term "sign"; but it seems to me convenient to make the triadic production of the interpretant essential to a "sign", calling the wider concept like a Jacquard loom, for example, a "quasi-sign". [...]
A sign may be defined as something (not necessarily existent) which is so determined by a second something called its Object that it will tend in its turn to determine a third something called its Interpretant in such a way that in respect to the accomplishment of some end consisting in an effect made upon the interpretant the action of sign is (more or less) equivalent to what that of the object might have been had the circumstances been different.
[...] How any sign, of whatsoever kind, mediates between an Object to some sort of conformity with which it is moulded, and by which it is thus determined, and an effect which the sign is intended to bring about and which it represents to be the outcome of the object influence upon it. It is of the first importance in such studies as these that the two correlates of the sign should be clearly distinguished : the Object by which the sign is determined and the Meaning, or as I usually call it, the Interpretant, which is determined by the sign, and through it by the object. The meaning may itself be a sign, a concept, for exemple, as may also the object. But everyboby who looks out of his eyes well knows that thoughts bring about tremendous physical effects, that are not, as such, signs. Feelings, too, may be excited by signs without thereby and theorein being themselves signs. We observe that the very same object may be several entirely different signs ; or in some way in other sign. [...] There are meanings that are feelings, meanings that are existent things or facts, and meanings that are concepts. [...]
[...] By a Sign, I mean anything that is, on the one hand, in some way determined by an object and, on the other hand, which determines some awareness, and this in such manner that the awareness is thus determined by that object. [...]
Of the distinction between the Objects, or better the "Originals" and the Interpretant of a Sign.
By "Sign" is meant any Ens which is determined by a single object or set of Objects called its Originals, all other than the Sign itself, and in its turn is capable of determining in a Mind something called its Interpretant, and that in such a way that the Mind is thereby mediately determined to some mode of conformity to the original or Set of originals. This is particularly intended to define (very imperfectly as yet) a complete Sign. But a complete sign has or may have Parts which partake of the nature of their whole; but often in a truncated fashion.
A Sign is in regard to its Interpretant in one or other of three grades of completeness, which may be called the Barely Overt, the Overter, and the Overtest. The Barely Overt of which a Name is an example does not expressly distinguish its original from its interpretant; nor its reference to either from the sign itself. The Overter sign of which an assertion is an exemple,... [phrase inachevée] Thus the Sign has a double function
1°/ to affect a mind which understands its "Grammar" or method of signification, which signification is its substance significate or Interpretant.
2°/ to indicate how to identify the conditions under which .... significate has the mode of being it is represented having [text unfinshed].
[...] Now any sign, of whatsoever kind, professes to mediate between an object, on the one hand, that to which it applies, and which is thus in a sense the cause of the sign, and, on the other hand, a Meaning, or to use a preferable technical term, an Interpretant, that which the sign expresses, the result which it produces in its capacity as sign. [...]
b - [...] Now any sign, of whatever kind, mediates between an object to some sort of conformity with which it is moulded, and which thus determines it, and an effect which it is intended to produce, and which it represents to be the outcome of the object. These two correlates of the sign have to be carefully distinguished. The former is called the object of the sign; the latter is the "meaning", or, as I usually term it, the "interpretant" of the sign. [...]
c - [...] Now the essential nature of a sign is that it mediates between its object which is supposed to determine it and to be, in some sense, the cause of it, and its meaning, or, as I prefer to say, in order to avoid certains ambiguities, its Interpretant which is determined by the sign, and is, in a sense, the effect of it; and which the sign represents to flow as an influences, from the object. [...]
d - [...] ...to which it is, therefore, conceived to be moulded, and by which to be determined, and an effect; on the other hand, which the sign is intended to bring about, representing it to be the outcome of the object influence upon it. I need not say that this influence is usually indirect and not of the nature of a force. [...]
e - [...] A sign is whatever there may be whose intent is to mediate between an utterer of it and an interpreter of it, both being repositories of thought, or quasi-minds, by conveying a meaning from the former to the latter. We may say that the sign is moulded to the meaning in the quasi-mind that utters it, where it was, virtually at least (i.e. if not in fact, yet the moulding of the sign took place as if it had been there) already an ingredient of thought.
But thought being itself a sign the meaning must have been conveyed to that quasi-mind, from some anterior utterer of the thought, of which the utterer of the moulded sign had been the interpreter. The meaning of the moulded sign being conveyed to its interpreter, became the meaning of a thought in that quasi-mind; and as these conveyed in a thought-sign required an interpreter, the interpreter of the moulded sign becoming the utterer of this new thought-sign".
f - I am now prepared to risk an attempt at defining a sign, -since in scientific inquiry, as in other enterprises, the maxim holds : nothing hazard, nothing gain. I will say that a sign is anything, of whatsoever mode of being, which mediates between an object and an interpretant; since it is both determined by the object relatively to the interpretant, and determining the interpretant in reference to the object, in such wise as to cause the interpretant to be determined by the object through the mediation of this "sign".
The object and the interpretant are thus merely the two correlates of the sign; the one being antecedent, the other consequent of the sign. Moreover, the sign being defined in terms of these correlative correlates, it is confidently to be expected that object and interpretant should precisely correspond, each to the other. In point of fact, we do find that the immediate object and emotional interpretant correspond, both being apprehensions, or are "subjective"; both, too, pertain to all signs without exception. The real object and energetic interpretant also correspond, both being real facts or things. But to our surprise, we find that the logical interpretant does not correspond with any kind of object. This defect of correspondance between object and interpretant must be rooted in the essential difference there is between the nature of an object and that of an interpretant; which difference is that former antecedes while the latter succeeds. The logical interpretant must, therefore, be in a relatively future tense.
[...] My idea of a sign has been so generalized that I have at length despaired of making anybody comprehend it, so that for the sake of being understood, I now limit it, so as to define a sign as anything which is on the one hand so determined (or specialized) by an object and on the other hand so determines the mind of an interpreter of it that the latter is thereby determined mediately, or indirectly, by that real object that determines the sign. Even this may well be thought an excessively generalized definition. The determination of the Interpreter's mind I term the Interpretant of the sign. [...]
It is clearly indispensable to start with an accurate and broad analysis of the nature of a Sign. I define a sign as an thing which is so determined by something else, called its Object, and so determines an effect upon a person, which effect I call its interpretant, that the latter is thereby mediately determined by the former. My insertion of "upon a person" is a sop to Cerberus, because I despair of making my own broader conception understood. [...]
A sign is a Cognizable that, on the one hand, is so determined (i.e. Specialized, bestimmt ), by something other than itself, called its object (or, in some cases, as if the Sign be the sentence "Cain killed Abel", in which Cain and Abel are equally Partial Objects, it may be more convenient to say that that which determines the Sign is the Complexus, or Totality, of Partial Objects. And in every case the object is accurately the Universe of which the special object is member, or part), while, on the other hand, it so determines some actual or potential Mind, the determination whereof I term the Interpretant created by the sign, that that interpreting mind is therein determined mediately by the Object.
A sign is a Cognizable that, on the one hand, is so determined (i.e. Specialized, bestimmt ), by something other than itself, called its object (or, in some cases, as if the Sign be the sentence "Cain killed Abel", in which Cain and Abel are equally Partial Objects, it may be more convenient to say that that which determines the Sign is the Complexus, or Totality, of Partial Objects. And in every case the object is accurately the Universe of which the special object is member, or part), while, on the other hand, it so determines some actual or potential Mind, the determination whereof I term the Interpretant created by the sign, that that interpreting mind is therein determined mediately by the Object.
Another endeavour to analyze a Sign.
A Sign is anything which represents something else (so far as it is complete) and if it represents itself it is as a part of another sign which represents something other than itself, and it represents itself in other circumstances, in other connections. A man may talk and he is a sign of that he relates, he may tell about himself as he was at another time. He cannot tell exactly what he is doing at that very moment. Yes, he may confess he is lying, but he must be a false sign, then. A sign, then, would seem to profess to represent something else.
Either a sign is to be defined as something which truly represents something or else as something which professes to represent something.
[...] I start by defining what I mean by a sign. It is something determined by something else its object and itself influencing some person in such a way that that person becomes thereby mediately influenced or determined in some respect by that Object.[...]
[...] Suffice it to say that a sign endeavours to represent, in part at least, an Object, which is therefore in a sense the cause, or determinant, of the sign even if the sign represents its object falsely. But to say that it represents its object implies that it affects a mind, and so affects it as, in some respect, to determine in that mind something that is mediately due to the Object.
That determination of which the immediate cause, or determinant, is the sign, and of which the mediate cause is the Object may be termed the Interpretant [...]
Signs, the only thing swith which a human being can, without derogation, consent to have any transaction, being a sign himself, are triadic; since a sign denotes a subject, and signifies a form of fact, which latter it brings into connexion with the former. [...]
Bya sign I mean anything whatever, real or fictive, which is capable of a sensible form, is applicable to something other than itself, that is already known, and that is capable of being so interpreted in another sign which I call its interpretant as to communicate something that may not have been previously known about its object there is thus a triadic relation between an sign, an Object, and an Interpretant.
55 - 1910 - C.P. 2-230 - Meaning .
The word sign will be used to denote an Object perceptible, or only imaginable, or even unimaginable in one sense -for the word "fast", which is a sign, is not imaginable, since it is not this word itself that can be set down on paper or pronounced, but only an instance of it, and since it is the very same word when it is written as it is when it is pronounced, but is one word when it means "rapidly" and quite another when it means "immovable", and a third when it refers to abstinence. But in order that anything should be a Sign, it must "represent" , as we say, something else, called its Object, although the condition that a sign must be other than its Object is perhaps arbitrary, since, if we insist upon it we must at least make an exception in the case of a sign that is a part of a sign. [...] A sign may have more than one Object.
Thus, the sentence "Cain killed Abel", which is a sign, refers at least as much to Abel as to Cain, even if it be not regarded as it should, as having "a killing" as a third object. But the set of objects may be regarded as making up one complex Object. In what follows and often elsewhere signs will be treated as having but one object each for the sake of dividing difficulties of the study. If a Sign is other than its object, there must exist, either in thought or in expression, some explanation or argument or other context, showing how -upon what system or for what reason the sign represents the Object or set of Objects that it does. Now the sign and the Explanation together make up another sign, and since the explanation will be a Sign, it will probably require an additional explanation, which taken together with the already enlarged Sign will make up a still larger sign; and proceeding in the same way, we shall, or should, ultimately reach a sign of itself, containing its own explanation and those of all its significant parts; and according to this explanation each such part has some other part as its Object. According to this every sign has, actually or virtually, what we may call a Precept of explanation according to which it is to be understood as a sort of emanation, so to speak, of its Object.
A logical Criticism of some articles of Religious Faith .
The word sign, as it will here be used, denotes any object of thought which excites any kind of mental action, whether voluntary or not, concerning something otherwise recognized. [...] Every sign denotes something, and the anything it denotes is termed an object of it. [... ] I term the idea or mental action that a sign exites and which it causes the interpreter to attribute to the Object or Objects of it, its interpretant. [...] For a Sign cannot denote an object not otherwise known to its interpreter, for the obvious reason that if he does not already know the Object at all, he cannot possess these ideas by means of which alone his attention can be narrowed to the very object denoted. Every object of experience excites an idea of some sort; but if that idea is not associated sufficiently and in the right way so with some previous experience so as to narrow the attention, it will not be a sign.
A Sign necessarily has for its Object some fragment of history, that is, of history of ideas. It must excite some idea. That idea may go wholly to narrowing the attention, as in such sign as "man", "virtue", "manner".
A Sketch of logical critic .
[...] In the first place, a "Representamen", like a word, -indeed, most words are representamens-, is not a single thing, but is of the nature of a mental habit, it consists in the fact that, something would be. The twenty odd "the" on ordinary page are all one and the same word, - that is, they are so many instances of a single word. Here are two instances of Representamens: "--killed--", "a man". The first of several characters which are each of them either essential to a sign being truly an instance of a Representamen or else necessary properties of such an object, is that it should have power to draw the attention of any mind that is fit to "interpret" it to two or more "Objects" of it. [ The first of the above examples of instances of representamens has four objects ; the second has two.] The second such character is that at least two of the objects must be other than the representamen. A closer examination than I have made would I am sure lead to a fuller description of the character. The third is the property that the interpreter of the representamen must have some collateral experiential acquaintance, direct or indirect, with each object of the Representamen before he can perform his function [...]
58 - v. 1911 - MS 676 : A Sketch of logical critics .
[...] If by a "sign" we mean anything of whatsoever nature that is apt to produce a special mental effect upon a mind in which certain associations have been produced -and I invariably use the word "association" as the original associationists did, for a mental habit, and never for the act or effect of associational suggestion when we must admit that a musical air and a command given to a soldier by his officer are signs, although it would seem that a logician is hardly otherwise concerned with such emotional and imperative signs, than that, as long as nobody else concerns himself with the analysis of the action of such signs, the logician is obliged to assume that office in order by the did of its contrast with the action of cognitional signs to perfect the definition of this latter. [...]
59 - 1911 - MS 854 - Notes on logical critique of the essential Articles of religious Faith (20.11.1911) .
Nature of a Sign . Its object is all that the sign recognize; since the sign cannot be understood until the Object is already identically known, though it may be indefinite. It so, it need only be known in its indefiniteness. The interpretant is the mental action on the Object that the sign excites.
For instance the word dog -meaning some dog, implies the knowledge that there is some dog, but it remains indefinite. The Interpretant is the somewhat indefinite idea of the characters that the "some dog" referred to has. And we have to distinguish between the Real Object and the Object as implied in the sign. The latter is some one of the dogs known already by direct experience or some one of the dogs which we more or less believe to exist.
The word dog does not excite any other notion than of the characters that ..... to possess.
The "Object" dog causes us to think of is such a dog as the person addressed has any notion of. But the real Object includes alternatively other dogs which are not known to the party addressed as yet but which he may come to know .
As to the characters we know it has four legs, is a carnivorous animal, etc.. and here we must distinguish then
- first the essential characters which the word implies -the essential interpretant.
- second the idea it actually does excite in the particular interpreter.
- third the characters it was intended specially to excite -perhaps only a part of the essential characters perhaps others not essential and which the word now excites though no such thing has hitherto been known.
In order to understand a Sign better we must consider that what it excites some sort of mental action about is in its Real Being either a history or a Part of a history and one part of it may be a Sign of another part.
Some Dog is a ....
Excites the idea of a Dog....is sign of a Dog and its Interpretant is forced by the interpreter own belief in the truth of the sign to regard its being a dog to admit that it is possible a ratter.
The sign may appeal to the Interpreter himself to assert that the Matter of Fact denoted does call for the....of certain character... or the Sign may exert a Force to cause the Interpreter to attach some Idea to the Object of the Sign.
60 - MS 670 :
A Sign, then, is anythin whatsoever -whether an Actual or a May-be or a Would-be,- which affects a mind, its Interpreter, and draw that interpreter's attention to some Object whether Actual, May-be or Would-be) which has already come within the sphere of his experience; and beside this purely selective action of a sign, it has a power of exciting the mind (whether directly by the image or the sound or indirectly) to some kind of feeling, or to effort of some kind or to thought; [...]
NON DATED TEXTS
The easiest of those which are of philosophical interest is the idea of a sign, or representation. A sign stands for something to the idea which it produces, or modifies. Or, it is a vehicle conveying into the mind something from without. That for which it stands is called its object; that which it conveys, its meaning; and the idea to which it gives rise, its interpretant. The object of representation can be nothing but a representation of which the first representation is the interpretant. But an endless series of representations, each representing the one behind it, may be conceived to have an absolute object at its limit. The meaning of a representation can be nothing but a representation. In fact, it is nothing but the representation itself conceived as stripped of irrelevant clothing. But this clothing never can be completely stripped off; it is only changed for something more diaphanous. So there is an infinite regression here. Finally, the interpretant is nothing but another representation to which the torch of truth is handed along; and as representation, it has its interpretant again. Lo, another infinite series.
62 - NEM IV - p. XXI - From MS.142.
A sign is a thing which is the representative, or deputy, of another thing for the purpose of affecting a mind [...]
63 - NEM IV - P. 239 - Kaina stoïcheia.
Any sign, B, which a sign, A, is fitted so to determine, without violation of its A's, purpose, that is, in accordance with the "Truth", even though it, B, denotes but a part of the objects of the sign, A, and signifies but a part of its, A's characters, I call an interpretant of A.
A sign is an object which stands for another to some mind. I propose to describe the characters of a sign. In the first place like any other thing it must have qualities which belong to it whether it be regarded as a sign or not thus a printed word is black, has a certain number of letters and those letters have certain shapes. Such characters of a sign I call its material quality. In the next place a sign must have some real connection with the thing it signifies so that when the object is present or is so as the sign signifies it to be the sign shall so signify it and otherwise not. [...] In the first place it is necessary for a sign to be a sign that it should be regarded as a sign for it is only a sign to that mind which so considers and if it is not a sign to any mind it is not a sign at all. It must be known to the mind first in its material qualities but also in its pure demonstrative application. That mind must conceive it to be connected with its object so that it is possible to reason from the sign to the thing. [...]
But at this point certain distinctions are called for. That which is communicated from the object through the Sign to the interpretant is a Form; that is to say, it is nothing like an existent, but is a power, is the fact that something would happen under certain conditions. This form is really embodied in the object, meaning that the conditional relation which constitutes the form is true of the form or it is in the Object. In the Sign it is embodied only in a representative sense, meaning that whether by virtue of some real modification of the Sign, or otherwise, the Sign becomes endowed with the power of communicating it to an interpretant. It may be in the interpretant directly, as it is in the Object, or it may be in the Interpretant dynamically, as behaviour of the Interpretant (this happens when a military officer uses the sign "Halt !" or "Forward march !" and his men simply obey him, perhaps automatically) or it may be in the Interpretant likewise only representatively. In existential graphs the Interpretant is affected in the last way; but for the present, it is best to consider only the common characters of all signs.
A Sign is an thing, A, which
(1) in addition to others characters of its own,
(2) stands in a dyadic relation Þ to a purely active correlate, B, and is also
(3) in a triadic relation to B for a purely passive correlate, C, this triadic relation being such as to determine C to be in a relation, µ, to B, the relation µ corresponding in a recognized way to the relation Þ, its dyadic relation to A would belong to it just the same even if A did not exist.
For instance, ...... the sign, the sentence "Let'songster of `Heliopolis' be our designation of the phenix" we may variously regard as B, either the phenix or the writer's determination, etc.. In any case howewer what is essential to the relation between the sentence and B is the writer's determination of mind to have the phenix called the songster of Heliopolis. This determination would be so shaped howewer whether expressed in this sentence or not. And the subsequent statement the sense in which certain correlates of a given relationship are said to be `active' or `passive' is that considering the different characters of all the correlates excepting only these that are immediately implied in the statement of the relationship none which involves only non-passives correlates will by immediate essential necessity vary with a variation of those involving only passive correlates; while no variation of which involve only non-active correlates will by immediate essential necessity carry with them variation of those which involve only active correlates; while by `active-passive' is meant active in respect to some correlates and passive in respect to others ........`active or passive' meaning........ active and ......without being active passive.
67 - MS 793 -[On Signs]1.
[...] which is communicated from the Object through the Sign to the Interpretant is a Form. It is not a singular thing; for if a Singular thing were first in the Object and afterward in the Interpretant outside the Object, it must thereby cease to be in the Object. The form that is communicated does not necessarily cease to be in one thing when it comes to be in a different thing, because its being is the being of a predicate. The Being of a Form consist in the truth of a conditional proposition. Under given circumstances something would be true. The Form is in the Object, one may say, entitatively, meaning that that conditional relation, or following of consequent upon reason, which constitutes the Form is literally true of the Object. In the Sign the Form may .... be embodied entitatively, but it must be embodied representatively, that is, in respect to the Form communicated, the Sign produce upon the interpretant an effect similar to that which the Object would under favorable circumstances.
68 - MS 793[On Signs] .
For the purpose of this inquiry a Sign may be defined as a Medium for the communication of a Form. It is not logically necessary that any thing possessing consciousness, that is, feeling or the peculiar commun quality of all our feeling should be concerned. But it is necessary that there should be two, if not three, quasi-minds, meaning things capable of varied determinations as to forms of the kind communicated.
As a medium the Sign is essentially in a triadic relation, to its Object which determines it and to its Interpretant which it determines. In its relation to the Object, the sign is passive, that is to say, its correspondence to the Object is brought about by on effect upon the sign, the Object remaining unaffected. On the other hand, in its relation to the Interpretant the sign is active determining the interpretant without bein itself thereby affected.
69 - MS 793 -[On Signs, quatre versions d'une certaine page 11] .
a - A Sign would be a Priman Secundan to something termed its Object and if anything were to be in a certain relation to the sign called being Interpretant to it, the Sign actively determines the Interpretant to be itself in a relation to the same Object, corresponding to its own.
b - b - A "Sign" is a genuinely genuine Tertian. It would generally be Priman in some characters, called its "Material Characters". But in addition, it is essentially (if only formally) Second to something termed its "Real Object", which is purely active in the Secundanity, being immediately unmodified by this secundanity; and these characters of the Real Object which are essential to the identity of the Sign constitute an ens rationis called the "Immediate Object". Moreover, the Sign is conceivably adapted to being Third to its Immediate Object for an ens rationis constituted thereby in the same (generic) relation to that Object in which the Sign itself stands to the same ; and this Third is termed the "Intended Interpretant", but the ... [unfinished]
c - A Sign would be in some respects Priman, and its determination as Priman are called its Material characters. But in addition it is Second to what is termed its Real Object, which is altogether active, and immediately unmodified by this Secundanity, and in so far as the Sign is second to it, it is termed the immediate Object. The Sign is conceivably adapted to being third to its Immediate Object for something in so far termed its Intended Interpretant; and the Sign only functions as such so far as the Intended Interpretant is Second to it for an Actual Interpretant which thus becomes adapted become a sign of the Immediate [there is a question mark above this word] Object for a further intended Interpretant, and in so far as the Interpretant is such Third it is termed Reflex Interpretant.
d - A "Sign" would be in some respects Priman, and its determinations as such are called its "Material characters". But in addition, it is Second to something termed its "Real Object", which is purely active being immediately unmodified by this Secundanity; and in so for as the sign is Second to it, it is termed the "Immediate Object" thereof. The Sign is conceivably adapted bo being Third to its Immediate Object for something which should thereby be brought into the generically same dyadic relation to that Object in which the Sign itself stands to that Object, and this Third is called the "Intended Interpretant"; but the Sign functions as such only in so far as the Intended Interpretant is Second to it and is Third to it for an existent termed the "Actual Interpretant", the modes of... [unfinished]
By a sign I mean any thing which is in any way, direct or indirect, so influenced by any thing (which I term its object) and which in turn influence a mind that this mind is thereby influenced by the Object; and I term that which is called forth in the mind the Interpretant of the sign. This explanation will suffice for the present; but distinctions will have to be drawn are long.
A mental representation is something which puts the mind into relation to an object. A representation generally (I am here defining my use of the term) is something which brings one thing into relation with another. The conception of third is here involved, and therefore, also, the conceptions of second or other and of first or an. A representation is in fact nothing but a something which has a third through an other. We may therefore consider an object :
1. as a something, with inward determinations;
2. as related to an other;
3. as bringing a second into relation to a third.
The most characteristic form of thirdness is that of a sign; and it is shown that every cognition is of the nature of a sign. Every sign has an object, which may be regarded either as it is immediately represented in the sign to be, and as it is in its own firstness. It is equally essential to the function of a sign that it should determine an Interpretant, or a second correlate related to the object of the sign as the sign is itself related to that object; and this interpretant may be regarded as the sign represents it to be, as it is in its pure secondness to the object, and as it is in its own firstness.
On the Classification of the Sciences .
A Representamen can be considered from three formal points of view, namely, first, as the substance of the representation, or the vehicle of the Meaning which is common to the three representamen of the triad, second, as the quasi-agent in the representation, conformity to which makes its Truth, that is, as the Natural Object, and third, as the quasi-patient in the representation, or that which modification in the representation make its Intelligence, and this may be called the Interpretant. Thus, in looking at a map, the map itself is the vehicle, the country represented is the Natural Object, and the idea excited in the mind is the Interpretant.
Furthermore, every representamen may be considered as a reagent, its intellectual character being neglected; and both representamen and reagent may be considered as quales, their relative character being neglected. This we do, for example, when we say that the word man has three letters.
by Robert Marty
ABSTRACT: We show that one can clearly distinguish two successive conceptions. The first (before 1905) which we qualify as "global triadic" and the second, more precise than the first, that we qualify as "analytic triadic" .
The 76 texts on the sign spread from1865 to 1911 (for 60 of them that are dated or whose dates are estimated). A brief study of the dispersion of the dated texts shows that more than 80% of them were produced after 1902, that is to say when Peirce was in his sixties. The production reaches a climax in l903, the year of the Lowell conferences. In addition, if one assesses them by their content, most of the non dated texts, and notably the eight definitions grouped in MS 793 are from the same period. Our purpose not being to study the evolution of Peirce's thinking in general, we will be interested only in the different conceptions of the sign that he proposes if, nevertheless, one can speak first of their differences before underlining their unity. Is it necessary to remind ourselves that the fundamental unity of these conceptions is confirmed by the constant reaffirmation of the triadic character of the sign? By describing the Peircean sign as triadic we simply highlight the presence in all implicit or explicit definition of the sign according to Peirce, three constitutive elements (but the datum of these three elements does not exhaust the Peircean concept of the sign since the relationship that links them together is lacking).
Peirce has varied the denomination of these three elements for reasons that he has sometimes clarified. We should not forget that he is the author of a very rigorous moral terminological (C.P.2-219 to 2 - 226). To designate the object of direct experience necessarily at the origin of all semiotic phenomena, Peirce uses the words "representation", "representamen" and especially "sign". He uses the term "representation" to this end only in texts n°1(1865), 6(1873) et 74 (n.d.), the other utilizations of the term designating the act or the fact of representing, as found in texts 10, 19, 27, 50, 52. In the text n°61, this word is given as a synonym for "sign ". There is thus no reason to retain this term.
On the other hand, it is interesting to examine very closely the different uses and distinctions between sign and representamen that Peirce first considers as synonymous (n° 9,1897) before making a distinction (n°19, l903) and finally deciding to abandon "representamen" (n°31, l905) since he explains that the popular usage of the word "sign" is very close to the exact sense of the scientific definition. In saying this, he makes the decision to put aside the formal distinction clearly established in l903 (n° 22, l903). We find the fundamental reason for abandoning this distinction in the statement, so often repeated by Peirce, that it was impossible to observe a single representamen that was not a sign. This conclusion is at odds with a number of authors, but in agreement, it seems to me, with Peirce (since from l905 he no longer uses the word representamen in any definitions except towards 1911 in the text n°57. However, the date attributed to this text being an estimation, it is possible to put it in doubt, and as in any event Peirce uses it in this text in a restricted sense, equivalent to legisign, there is no need to preserve this " horrible word " and "sign" should be quite suitable. There would have perhaps been some interest, on the other hand, in preserving representamen so as to concretize the different conceptualizations of semiotic phenomena as between the Saussuro-hjelmslevian tradition and the Peircean tradition. But the adoption of this viewpoint would be a sort of renunciation of the debates on the profound nature of these phenomena according to these two traditions ; the passive acceptance of the fact that both traditions should develop independently would thus deprive us of the clarity that the opening of conflicts can bring about in the semiotic field.
To designate the object of the sign, Peirce employs on nearly every occasion the word "object" accompanied with considerations that render it, explicitly or implicitly, that which is connected to this object of direct experience that is the sign. Sometimes Peirce designates it by the expression "some thing " and even in the text n°23 the sign is said to represent an aspect of the "True" (the "Truth", the true universe), another representamen in the text n°21, and a subject in the n°53. Moreover, the object is often qualified: Real, Natural or Original in addition to the distinction between immediate object and dynamic objects.
Despite these remarks, there is no problem in denominating as "Object" this other object whose presence to the mind produced by the perception of the sign is characteristic of semiotic phenomena. It is clear that a third element is needed because it is essential in semiotic phenomena to define an element capable of explaining the necessary connection of the two objects that are potentially present to the mind (the perceived sign, as such, and the object to which it is connected). For if the sign, an object of direct experience, is distinguishable because it evokes another object different from itself, because it enables a supplementary perceptive choice (at least) it constitutes, by this very fact, an association between these two objects. That the sign is one of the two objects does not change anything in the matter; it both exists for itself and exists for another. But this association can be conceived only in the mind and by the mind to which the two objects are present. In a sign in actu this association is truly a matter of fact; it is a psychic fact that the mind that constructs two different perceptual judgements on the same percept is in a special state, different from that which it is in, in the case of ordinary phenomena, that is to say in the simple presentation of an object, due to the fact of this dual presence (it is the thesis that I develop in my work in French "The Algebra of signs" (1990, John Benjamins, Amsterdam/Philadelphia). One can say that this special state of mind gives at this very instant a real existence to this association, even in the most "natural " cases. Friday's footprint in the sand stands for a human presence only because of the association in Robinson's mind, even if its production and therefore its existence are totally independent of his mind. In every sign there intervenes therefore the determination of a mind, distinct from the two objects, which is therefore an element necessarily implied in the factuality of the sign and without which one cannot hope to describe semiotic phenomena correctly. The subject is therefore implied in a certain manner in this approach. It is necessary therefore to attach a third element to the Sign and to the Object. Peirce gives it the name of Interpretant. Now, let us examine the various denominations by which he himself grasped this necessity. Note immediately that the last sentence of the textn°6 (1873) covers exactly the argument we have just developed: "The idea of the representation itself excites in the mind another idea and in order that it may do this it is necessary that some principle of association between the two ideas should already be established in that mind". One finds again this idea in the text n°64 (n.d) :
That mind must conceive it to be connected with its object so that it is possible to reason from the sign to the thing.
and in the text n° 58 (v.1911), the Interpretant is
"a special mental effect upon a mind in which certain associations have been produced".
A systematic list of the words that Peirce uses to give content to the concept of the interpretant shows that he attributes the following characteristics, according to what he is saying at that moment and to the maturation of his thinking: - it is a thought or interpretant thought in texts n°, 8, l0, 18, 28.
-- it is a determination of a mind or quasi -mind or an influence on a person or a mind, this determination or influence being realised through the sign, the object by being the mediate cause in texts n°34, 37, 40(a,b,c,e,f,) 52.
- It is a Third that according to the case is a third correlate of a triadic relationship or a "Tertian", (that is to say a member of the Third universe, a Thirdness) in texts n°13, 15, 20, 22, 36, 69(b,c,d,e). Moreover, in n°30 iit is described as a "passive" correlate and in n°76 it is a quasi - patient.
One sees that these characteristics (by excluding the last that is of a radically different nature), can be classified in two groups:
- Those that refer to a sign in actu, that describe therefore this third element of the semiotic phenomenon in its particularity and that are practically reducible to an effect on a person or again to a determination of a mind, in the here and now of the perception.
- Those that refer to an abstract sign which come from the logical analysis of the phenomenon and form part of a formal construction, in which the interpretant is described as a correlate of a triadic relationship.
Peirce had a great deal of difficulty getting people to accept this conception, quite banal today, of a formal model of the sign. In his letter to Lady Welby dated 23 December l908 he complains about his difficulties by writing : "I have added 'on a person' so as to throw a cake to Cerberus, because I despair of making people understand my own conception which is larger. " In conclusion then, peircean conceptions of the sign lead us to retain three fundamental elements as theoretical universals resulting from the logical analysis of semiotic phenomena,:
- the Sign S, an object of direct experience ("external"or "internal"object").
- the Object O, present in the semiotic phenomenon because it is connected with the sign.
- the Interpretant I, present because it is a mental element which ensures this connection.
The reader will have noticed that these groups and subgroups of definitions possess common elements since the characteristics fundamental to them are not exclusive of each other. However by observing the placings of texts constituting the subgroups and by reminding ourselves that numbers l to 60 are classified by chronological order, this distribution shows a significant change, if not of doctrine, at least of his approach to this connection of the Sign to its Object. It suffices indeed to observe the pre-eminence from n°29 (l905) of the characterization of this connection in terms of determination of S by O to arrive at the conclusion that Peirce has decided to take into account, round about l905, the dissymetric character of this relationship, which he has expressed by writing that if, in a sign, O acts on S, the reverse is not necessarily true. The consequence of this change will be the abandoning of the central position granted to the triad in the global approach to the sign. Indeed, to define a priori the sign as triadic implies that diadic relationships between two elements that are induced by the triadic relationship are symmetrical. Therefore if one wants to preserve the dissymmetry of this relationship, it is necessary either to abandon the idea of basing the sign on the notion of triad, or to add correctives (which would be difficult), or to change the perspective, which does not imply the renunciation of the triadicity but simply causes it to intervene at another level. We will see later, that a third approach, based on the notion of communication, will tend to unify the two precedent perspectives.
With regard to the connection between Sign and Interpretant, it is invariably conceived, each time than it is evoked, as a relationship of determination (when it is evoked in a formal model), an effect on interpreter or a determination of the mind of an interpreter (when it concerns the description of a sign in actu). In text n°49 (l909) lthe Interpretant is even called "a creature of the sign ", a conception which is problematic if one thinks of the necessity, many time underlined by Peirce, that an association is a prerequisite in the mind in order that a sign might function as such, which obviously excludes the possibility that the sign could create the Interpretant ex-nihilo. It is what Peirce resumed in this text by specifying that the Interpretant is created by the Sign "in its capacity to support its determination by the Object". In one of his most formal approaches in which the triadic relationship is his point of departure, (C.P.2-233 and s.q.q., Division of Triadic Relations) Peirce defines the Representamen ( see n°22, v.l903) as the first correlate of an authentic triadic relationship. In other words he considers that this first correlate determines the third correlate. Thus, at that moment, he approaches the sign through its triadicity to which he adds a corrective: the determination of the interpretant by the sign, the connections signs-object and object-interpretant being induced by the triadic relationship, the sign itself being a particular representamen , namely a representamen that determines a particular interpretant that is the "act of cognition of a mind".
By taking into account the dissymmetry of the relationship Object-Sign he has therefore, as we have noticed, abandoned the triadicity as founder principle and has resorted to a new notion, linked to the higthlighting of successive determinations (of the Sign by the Object and the Interpretant by the Sign) in the analysis of the sign in actu, the notion of mediation. It concerns a resumption of this notion already present in 1867 (" mediate representation "in text n°8) and in l902 (" the authentic mediation is the character of a sign"in text n°13).In l904, triadicity and mediation appear in the same text (n°28). However one can observe that in the majority of texts after l905 that mention the two determinations cited above, one of words "mediation"or"medium "or the verb"to mediate"is present (texts n°: 33, 37, 39, 40 (a,b,c,e,f) 46, 47, 48, 49,51, 52). It concerns a new theoretical approach (because the term" triad"does not figure in any of these texts) that is based this time on the determination of the Interpretant by the Object through the sign. This conception is partially clarified and formalized in text n° 30 (l905) in which the triad is still present: the Sign is presented there as a passive correlate in its relationship to the Object, which relationship is incorporated in a triadic relationship in such a manner that the Interpretant is put in a diadic relationship with the Object, induced by this triadic relationship. What does not figure in this definition assuredly the most formalized of all, (and that one finds in text n°66 in the undated manuscript n°793) is precisely the determination of the Interpretant by the Sign.
It is in text n°32 that the change appears fully with the notion of "medium of communication ". The next text (n°33) is more precise: a sign is"a medium for the communication or the extension of a form (or figure)". One finds this idea of form in n° 53 and 54 (19l0). It would seem that Peirce has attempted to explain the fact that the determination of the Sign by the Object was such that it produced the indirect determination of the Interpretant by the Object taking into account that a certain "form"was present in each of the three elements of the sign, as soon as the sign was established, and that the process of establishment of the sign consisted in communicating (or conveying) this form from the Object to the Interpreter through the Sign. This step does not exclude triadicity insofar as it is precisely the presence of this "form"that, we think, allows us to link triadically the three elements of the semiotic phenomenon (by being incorporated in each of them). It would be the ground evoked by Peirce in 2-228 (text n°9, v.1897). One sees therefore that the two main theoretical approachs that we have just elucidated in this group of texts, are not exclusive. In conclusion we will distinguish therefore, without opposing them, two Peircean conceptions of the sign: - a conception that, for convenience, we will call "global triadic " derived from an analysis of semiotic phenomena which considers as essential the fact that the three elements therein are necessarily linked by a triadic relationship. - a conception that we will call" analytic triadic" derived from a finer analysis in terms of the determination of some elements by others (of the sign by the object and the interpretant by the sign), the interplay of these two determinations leading to the establishment of a triadic relationship between the three elements necessarily present in semiotic phenomena (it is the presence of the conveyed Form in the course of these successive determination that creates the triadic relationship). To grasp this second conception better it is necessary to clarify what Peirce understands by "determination" in the precise case of the sign, or, in view of the difficulty of the task, to try to discern this notion better. Because the explanations given by Peirce as to the sense in which he uses the words "active" and "passive"in texts n°30 and 66 appear to us to be no longer operative. To the extent that we have been able to understand his thinking, it seems to us that Peirce considers that there is character determination of one correlate by those of another, the correlate B being active with reference to the correlate A, if all characters of this latter which are involved in the semiotic phenomenon are implied by the characters of B. Friday's footprint in the sand perfectly illustrates this notion since it is just what it is, that is to say possessed of characters that make it a sign, because the foot that has produced it has communicated them without being modified itself, and it is thus a purely active correlate. The imprint itself is a purely passive correlate for opposite reasons . However if now one photographs this imprint, it is going to produce an image on the film which owes all its characters to the imprint itself. In relation to this photographic image the imprint will be therefore an active correlate and it is clear that, for Peirce, the interpretant C is a purely passive correlate determined by the imprint, this interpretant, triadic in nature, being such that it incorporates, as an induced diadic relationship, the diadic relationship established between Friday's foot and its imprint. However the example that we quote is particular, it is a scholastic example. Nevertheless it is, we think, by generalizing the case of signs of this type ( index) that Peirce obtained the definition n° 30.In others texts he has used terms that allow us to higthlight somewhat this conception:
-in 46 and 48 the sign is said to be specialized (that Peirce strengthens by calling it in German "bestimmt") and in 47 and 48 he writes that the determination of the Sign by the Object is such that consequently it determines the Interpretant, what means that if the Sign is passive in relation to the Object and active when related to the Interpretant, it owes this last possibility to the action of the Object, as a pool ball becomes capable of moving another after having itself been knocked by another one. Moreover, in text n°65, Peirce makes it clear that when the Form which comes from the Object is incorporated in the Sign the former becomes "endowed with the power to communicate it to an Interpretant".
- but it is certainly in text n° 40 f, that Peirce clearly presents as an attempt to define the sign, that his conception of determination in semiotic phenomena is best expressed while being probably the most difficult to formalize: the sign, he writes, is both " determined by the object with respect to the interpretant and it determines the interpretant in reference to the object, in such a manner that the interpretant is determined by the object as a cause through the mediation of this sign". One sees that the determinations of elements one by other (of the Sign by the Object, the Interpretant by the Sign) are constructed dependent on the third, for lack of which, the semiotic phenomenon would be reduced to the compositionof two independent successive determinations, in contradiction with the consistant doctrine of Peirce. It is by taking into account the whole of Peirce's contributions, of which it is unnecessary to underline the wealth, the power but also the difficulty, that we have taken the risk (since we are agreement with Peirce -cf 40 f - that in scientific matters, as in other enterprises the maxim:"no risk, no profit" is valid ") of putting forward a formal definition of the sign that is operative and also mathematically formalizable, to reach as far as possible towards authentically scientific semiotics (cf " The Algebra of Signs " an essay in scientific semiotics according to Charles Sanders Peirce) |
|— Capital city —|
|Founder||Abu Jafar al-Mansur|
|• Body||Baghdad City Advisory Council|
|• Mayor||Saber Nabet Al-Essawi|
|• Total||204.2 km2 (78.8 sq mi)|
|Elevation||34 m (112 ft)|
|• Estimate (2011)||7,216,040|
|Time zone||Arabia Standard Time (UTC+3)|
|• Summer (DST)||No DST (UTC)|
|Postal code||10001 to 10090|
Baghdad (Arabic: بغداد, Baġdād, IPA: [bæɣˈdæːd]) is the capital of the Republic of Iraq, as well as the coterminous Baghdad Province. The population of Baghdad as of 2011 is approximately 7,216,040, making it the largest city in Iraq, the second largest city in the Arab World (after Cairo, Egypt), and the second largest city in Western Asia (after Tehran, Iran). According to the government, which is preparing for a census, the population of the country has reached 35 million, with 9 million in the capital.
Located along the Tigris River, the city was founded in the 8th century and became the capital of the Abbasid Caliphate. Within a short time of its inception, Baghdad evolved into a significant cultural, commercial, and intellectual center for the Islamic World. This in addition to housing several key academic institutions (e.g. House of Wisdom) garnered the city a worldwide reputation as the "Center of Learning". Throughout the High Middle Ages, Baghdad was considered to be the largest city in the world with an estimated population of 1,200,000 people. The city was largely destroyed at the hands of the Mongol Empire in 1258, resulting in a decline that would linger through many centuries due to frequent plagues and multiple successive empires. With the recognition of Iraq as an independent state (formerly the British Mandate of Mesopotamia) in 1938, Baghdad gradually regained some of its former prominence as a significant center of Arabic culture.
In contemporary times the city has often faced severe infrastructural damage, most recently due to the American-led foreign occupation in March 2003 that lasted until December 2011 and the subsequent sectarian violence. In recent years the city has been frequently subjected to insurgency activities and terrorist attacks. Though the nation continues to work toward rebuilding and reconciliation, as of 2012 Baghdad continues to be listed as one of the least hospitable places in the world to live and was ranked by Mercer as the worst of 221 major cities as measured by quality-of-life.
The name Baghdad is pre-Islamic and its origins are under some dispute. Some say it comes from an Aramaic phrase that means "sheep enclosure". Others contend that the word comes from ancient Persian: "bagh" meaning God, and "dad" meaning gift. "The gift of God...." During at least one point in history, it certainly seemed so. The name has been found as Baghdadu on Assyrian cuneiform records of the 9th century BC, and on Babylonian bricks bearing the Royal Seal of King Nebuchadnezzar (6th century BC).[better source needed]. However it is related to previous settlements, which did not have any political or commercial power.
When the Abbasid caliph al-Mansur founded a completely new city for his capital, he chose the name Madinat al-Salaam or City of Peace. This was the official name on coins, weights, and other official usage, although the common people continued to use the old name. By the 11th century, "Baghdad" became almost the exclusive name for the world-renowned metropolis.
There are several rival theories as to the etymology of the specific name Baghdad. The most widely accepted among these is that the name is a Middle Persian compound of Bag "god" + dād "given", translating to "God-given" or "God's gift", from which comes Modern Persian Baɣdād. This in turn can be traced to Old Persian. A less probable guess is the Persian compound Bāğ "garden" + dād "fair", translating to "The fair garden".
After the fall of the Persian Sassanid empire, the victorious Muslim rulers wanted their own capital. Choosing a site north of the Persian capital of Ctesiphone, on 30, July 734 the caliph Al Mansur commissioned the construction of the city and it was built under the supervision of the Barmakids. Mansur believed that Baghdad was the perfect city to be the capital of the Islamic empire under the Abbasids. Mansur loved the site so much he is quoted saying, "This is indeed the city that I am to found, where I am to live, and where my descendants will reign afterward".
The city's growth was helped by its location, which gave it control over strategic and trading routes, along the Tigris. A reason why Baghdad provided an excellent location was the abundance of water and the dry climate. Water exists on both north and south ends of the city gates, allowing all households to have a plentiful supply, which was very uncommon during this time.
Baghdad eclipsed Ctesiphon, the capital of the Persian Empire, which was located some 30 km (19 mi) to the southeast. Today, all that remains of Ctesiphon is the shrine town of Salman Pak, just to the south of Greater Baghdad. Ctesiphon itself had replaced and absorbed Seleucia, the first capital of the Seleucid Empire. Seleucia had earlier replaced the city of Babylon.
In its early years the city was known as a deliberate reminder of an expression in the Qur'an, when it refers to Paradise. It took four years to build (734-738). Mansur assembled engineers, surveyors, and art constructionists from around the world to come together and draw up plans for the city. Over 100,000 construction workers came to survey the plans; many were distributed salaries to start the building of the city. July was chosen as the starting time because two astrologers, Naubakht Ahvaz and Mashallah, believed that the city should be built under the sign of the lion, Leo. Leo is associated with fire and symbolises productivity, pride, and expansion.
The bricks used to make the city were 18 inches (460 mm) on all four sides. Abu Hanifa was the counter of the bricks and he developed a canal, which brought water to the work site for the use of both human consumption and the manufacturing of the bricks. Marble was also used to make buildings throughout the city, and marble steps led down to the river's edge.
The basic framework of the city consists of two large semicircles about 19 km (12 mi) in diameter. The city was designed as a circle about 2 km in diameter, leading it to be known as the "Round City". The original design shows as single ring of residential and commercial structures along the inside of the city walls, but the final construction added another ring inside the first. Within the city there were many parks, gardens, villas, and promenades. In the center of the city lay the mosque, as well as headquarters for guards. The purpose or use of the remaining space in the center is unknown. The circular design of the city was a direct reflection of the traditional Persian Sasanian urban design. The Sasanian city of Gur in Fars, built 500 years before Baghdad, is nearly identical in its general circular design, radiating avenues, and the government buildings and temples at the centre of the city. This style of urban planning contrasted with Ancient Greek and Roman urban planning, in which cities are designed as squares or rectangles with streets intersecting each other at right angles.
The surrounding walls
The four surrounding walls of Baghdad were named Kufa, Basra, Khurasan, and Damascus; named because their gates pointed in the directions of these destinations. The distance between these gates was a little less than 1.5 miles (2.4 km). Each gate had double doors that were made of iron; the doors were so heavy it took several men to open and close them. The wall itself was about 44 m thick at the base and about 12 m thick at the top. Also, the wall was 30 m high, which included merlons, a solid part of an embattled parapet usually pierced by embrasures. This wall was surrounded by another wall with a thickness of 50 m. The second wall had towers and rounded merlons, which surrounded the towers. This outer wall was protected by solid glacis, which is made out of bricks and quicklime. Beyond the outer wall was a water-filled moat.
Golden Gate Palace
In the middle of Baghdad, in the central square was the Golden Gate Palace. The Palace was the residence of the caliph and his family. In the central part of the building was a green dome that was 39 m high. Surrounding the palace was an esplanade, a waterside building, in which only the caliph could come riding on horseback. In addition, the palace was near other mansions and officer's residences. Near the Gate of Syria a building served as the home for the guards. It was made of brick and marble. The palace governor lived in the latter part of the building and the commander of the guards in the front. In 813, after the death of caliph Al-Amin the palace was no longer used as the home for the caliph and his family. The roundness points to the fact that it was based on Arabic script. The two designers who were hired by Al-Mansur to plan the city's design were Naubakht, a Zoroastrian who also determined that the date of the foundation of the city would be astrologically auspicious, and Mashallah, a Jew from Khorasan, Iran.
The Abbasids and the round city
The Abbasid Caliphate was based on their being the descendants of the uncle of Muhammad and being part of the Quraysh tribe. They used Shi'a resentment, Khorasanian movement, and appeals to the ambitions and traditions of the newly conquered Persian aristocracy to overthrow the Umayyads. The Abbasids sought to combine the hegemony of the Arabic tribes with the imperial, court, ceremonial, and administrative structures of the Persians. The Abbasids considered themselves the inherittures and the need of Mansur to place the capital in a place that was representative of Arab-Islamic identity by building the House of Wisdom, where ancient texts were translated from their original language, such as Greek, to Arabic. Mansur is credited with the "Translation Movement" for this. Further, Baghdad is also near the ancient Sassanid imperial seat of Ctesiphon on the Tigris River.
A centre of learning (8th to 13th centuries)
Within a generation of its founding, Baghdad became a hub of learning and commerce. The House of Wisdom was an establishment dedicated to the translation of Greek, Middle Persian and Syriac works. Scholars headed to Baghdad from all over the Abbasid empire, facilitating the introduction of Persian, Greek and Indian science into the Arabic and Islamic world at that time. Baghdad was likely the largest city in the world from shortly after its foundation until the 930s, when it was tied by Córdoba. Several estimates suggest that the city contained over a million inhabitants at its peak. Many of the One Thousand and One Nights tales are set in Baghdad during this period.
The end of the Abbasids in Baghdad
By the 10th century, the city's population was between 1.2 million and 2 million. Baghdad's early meteoric growth eventually slowed due to troubles within the Caliphate, including relocations of the capital to Samarra (during 808–819 and 836–892), the loss of the western and easternmost provinces, and periods of political domination by the Iranian Buwayhids (945–1055) and Seljuk Turks (1055–1135).
The Seljuks were a clan of the Oghuz Turks from the Central Asia that converted to the Sunni branch of Islam. In 1040, they destroyed the Ghaznavids, taking over their land and in 1055, Tughril Beg, the leader of the Seljuks, took over Baghdad. The Seljuks expelled the Buyids dynasty of Shiites that ruled for some time and took over power and control of Baghdad. They ruled as Sultans in the name of the Abbasid caliphs (they saw themselves as being part of the Abbasid regime) Tughril Beg saw himself as the protector of the Abbasid Caliphs.
In 1058, Baghdad was captured by the Fatimids under the Turkish general Abu'l-Ḥārith Arslān al-Basasiri, an adherent of the Ismailis along with the 'Uqaylid Quraysh. Not long before the arrival of the Saljuqs in Baghdad, al-Basasiri petitioned to the Fatimid Imam-Caliph al-Mustansir to support him in conquering Baghdad on the Ismaili Imam's behalf. It has recently come to light that the famed Fatimid da'i al-Mu'ayyad al-Shirazi had a direct role in supporting al-Basasiri and helped the general to succeed in taking Mawṣil, Wāsit and Kufa. Soonafter, By December 1058, a Shi'i adhān (call to prayer) was implemented in Baghdad and a khutba (sermon) was delivered in the name of the Fatimid Imam-Caliph. Despite his Shi'i inclinations, Al-Basasiri received support from Sunnis and Shi'is alike, for whom opposition to the Saljuq power was a common factor.
On February 10, 1258, Baghdad was captured by the Mongols led by Hulegu, a grandson of Chingiz Khan (Genghis Khan) during the sack of Baghdad. Many quarters were ruined by fire, siege, or looting. The Mongols massacred most of the city's inhabitants, including the caliph Al-Musta'sim, and destroyed large sections of the city. The canals and dykes forming the city's irrigation system were also destroyed. The sack of Baghdad put an end to the Abbasid Caliphate, a blow from which the Islamic civilization never fully recovered.
At this point Baghdad was ruled by the Il-Khanids, the Mongol rulers of Iran. In 1401, Baghdad was again sacked, by the Central Asian Turkic conqueror Timur ("Tamerlane"). When his forces took Baghdad, he spared almost no one, and ordered that each of his soldiers bring back two severed human heads. It became a provincial capital controlled by the Mongol Jalayirid (1400–1411), Turkic Kara Koyunlu (1411–1469), Turkic Ak Koyunlu (1469–1508), and the Iranian Safavid (1508–1534) dynasties.
Ottoman era (16th to 19th centuries)
In 1534, Baghdad was captured by the Ottoman Turks. Under the Ottomans, Baghdad fell into a period of decline, partially as a result of the enmity between its rulers and Iranian Safavids, which did not accept the Sunni control of the city. Between 1623 and 1638, it returned to Iranian rule before falling back into Ottoman hands.
For a time, Baghdad had been the largest city in the Middle East. The city saw relative revival in the latter part of the 18th century under a Mamluk government. Direct Ottoman rule was reimposed by Ali Ridha Pasha in 1831. From 1851-1852 and from 1861–1867, Baghdad was governed, under the Ottoman Empire by Mehmed Namık Pasha. The Nuttall Encyclopedia reports the 1907 population of Baghdad as 185,000. Aside from ethnically Arab Iraqis, the city was also home to a substantial ancient Jewish community, which comprised over a quarter of the city's population (this proportion would grow in later years).
20th century
Baghdad and southern Iraq remained under Ottoman rule until 1917, when captured by the British during World War I. From 1920, Baghdad became the capital of the British Mandate of Mesopotamia and, after 1932, Baghdad was the capital of the Kingdom of Iraq. Iraq was given formal independence in 1932 and increased autonomy in 1946. The city's population grew from an estimated 145,000 in 1900 to 580,000 in 1950.
On 1 April 1941 members of the "Golden Square" and Rashid Ali staged a coup in Baghdad. Rashid Ali installed a pro-German and pro-Italian government to replace the pro-British government of Regent Abdul Ilah. On 31 May, after the resulting Anglo-Iraqi War and after Rashid Ali and his government had fled, the Mayor of Baghdad surrendered to British and Commonwealth forces.
On 14 July 1958, members of the Iraqi Army under Abdul Karim Kassem staged a coup to topple the Kingdom of Iraq. King Faisal II, former Prime Minister Nuri al-Said, former Regent Prince Abdul Ilah, members of the royal family, and others were brutally killed during the coup. Many of the victim's bodies were then dragged through the streets of Baghdad.
During the 1970s, Baghdad experienced a period of prosperity and growth because of a sharp increase in the price of petroleum, Iraq's main export. New infrastructure including modern sewerage, water, and highway facilities were built during this period. However, the Iran–Iraq War of the 1980s was a difficult time for the city, as money was diverted by Saddam Hussein to the army and thousands of residents were killed. Iran launched a number of missile attacks against Baghdad in retaliation for Saddam Hussein's continuous bombardments of Tehran's residential districts.
In 1991 and 2003, the Gulf War and the 2003 invasion of Iraq caused significant damage to Baghdad's transportation, power, and sanitary infrastructure as the US-led coalition forces launched massive aerial assaults in the city in the two wars.
Main sights
Points of interest include the National Museum of Iraq whose priceless collection of artifacts was looted during the 2003 invasion, and the iconic Hands of Victory arches. Multiple Iraqi parties are in discussions as to whether the arches should remain as historical monuments or be dismantled. Thousands of ancient manuscripts in the National Library were destroyed wh under Saddam's command.
Baghdad Zoo
The Baghdad Zoo was the largest zoo in the Middle East. Within eight days following the 2003 invasion, however, only 35 of the 650 animals in the facility survived. This was a result of theft of some animals for human food, and starvation of caged animals that had no food invasion, South African Lawrence Anthony and some of the zoo keepers cared for the animals and fed the carnivores with donkeys they had bought locally. Eventually, L. Paul Bremer, Director of the Coalition Provisional Authority in Iraq from May 11, 2003 to June 28, 2004 ordered protection of the zoo and U.S. engineers helped to reopen the facility.
Geography and climate
The city is located on a vast plain bisected by the River Tigris. The Tigris splits Baghdad in half, with the eastern half being called 'Risafa' and the Western half known as 'Karkh'. The land on which the city is built is almost entirely flat and low-lying, being of alluvial origin due to the periodic large floods which have occurred on the river.
Baghdad has a subtropical arid climate (Köppen climate classification BWh) and is, in terms of maximum temperatures, one of the hottest cities in the world. In the summer from June to August, the average maximum temperature is as high as 44 °C (111 °F) accompanied by blazing sunshine: rainfall has in fact been recorded on fewer than half a dozen occasions at this time of year and has never exceeded 1 millimetre (0.04 in). Temperatures exceeding 50 °C (122 °F) in the shade are by no means unheard of, and even at night temperatures in summer are seldom below 24 °C (75 °F). Because the humidity is very low (usually under 10%) due to Baghdad's distance from the marshy Persian Gulf, dust storms from the deserts to the west are a normal occurrence during the summer.
Winters boast mild days and variable nights. From December to February, Baghdad has maximum temperatures averaging 15.5 to 18.5 °C (60 to 65 °F), though highs above 70 °F (21 °C) are not unheard of. Morning temperatures can be chilly: the average January low is 3.8 °C (38.8 °F) but lows below freezing only occur a couple of times per year.
Annual rainfall, almost entirely confined to the period from November to March, averages around 150 mm (5.91 in), but has been as high as 338 mm (13.31 in) and as low as 37 mm (1.46 in). On January 11, 2008, light snow fell across Baghdad for the first time in memory.
|Climate data for Baghdad|
|Average high °C (°F)||16
|Average low °C (°F)||3.8
|Rainfall mm (inches)||26
|Avg. rainy days (≥ 0.1 mm)||5||6||4||2||0||0||0||0||0||1||5||6||29|
|Mean monthly sunshine hours||192.2||203.3||244.9||255.0||300.7||348.0||347.2||353.4||315.0||272.8||213.0||195.3||3,240.8|
|Source #1: Climate & Temperature|
|Source #2: World Meteorological Organisation (UN)|
Administrative divisions
The city of Baghdad has 89 official neighbourhoods within 9 districts. These official subdivisions of the city served as administrative centres for the delivery of municipal services but until 2003 had no political function. Beginning in April 2003, the U.S. controlled Coalition Provisional Authority (CPA) began the process of creating new functions for these. The process initially focused on the election of neighbourhood councils in the official neighbourhoods, elected by neighbourhood caucuses.
The CPA convened a series of meetings in each neighbourhood to explain local government, to describe the caucus election process and to encourage participants to spread the word and bring friends, relatives and neighbours to subsequent meetings. Each neighbourhood process ultimately ended with a final meeting where candidates for the new neighbourhood councils identified themselves and asked their neighbours to vote for them.
Once all 88 (later increased to 89) neighbourhood councils were in place, each neighbourhood council elected representatives from among their members to serve on one of the city's nine district councils. The number of neighbourhood representatives on a district council is based upon the neighbourhood's population. The next step was to have each of the nine district councils elect representatives from their membership to serve on the 37 member Baghdad City Council. This three tier system of local government connected the people of Baghdad to the central government through their representatives from the neighbourhood, through the district, and up to the city council.
The same process was used to provide representative councils for the other communities in Baghdad Province outside of the city itself. There, local councils were elected from 20 neighbourhoods (Nahia) and these councils elected representatives from their members to serve on six district councils (Qada). As within the city, the district councils then elected representatives from among their members to serve on the 35 member Baghdad Regional Council.
The first step in the establishment of the system of local government for Baghdad Province was the election of the Baghdad Provincial Council. As before, the representatives to the Provincial Council were elected by their peers from the lower councils in numbers proportional to the population of the districts they represent. The 41 member Provincial Council took office in February, 2004 and served until national elections held in January 2005, when a new Provincial Council was elected.
This system of 127 separate councils may seem overly cumbersome but Baghdad Province is home to approximately seven million people. At the lowest level, the neighbourhood councils, each council represents an average of 75,000 people.
The nine District Advisory Councils (DAC) are as follows:
- Sadr City (Thawra)
- Al Rashid
- New Baghdad (Tisaa Nissan) (9 April)
The nine districts are subdivided into 89 smaller neighborhoods which may make up sectors of any of the districts above. The following is a selection (rather than a complete list) of these neighborhoods:
- Hurriya City
- Bab Al-Moatham
- Hayy Ur
- Hayy Al-Jami'a
- Al Khadhraa
- Hayy Al-Jihad
- Hayy Al-A'amel
- Hayy Aoor
- Hayy Al-Shurtta
- Jesr Diyala
- Abu Disher
- Raghiba Khatoun
- Arab Jijur
|This section requires expansion. (December 2009)|
Reconstruction efforts
Most Iraqi reconstruction efforts have been devoted to the restoration and repair of badly damaged urban infrastructure. More visible efforts at reconstruction through private development, like architect and urban designer Hisham N. Ashkouri's Baghdad Renaissance Plan and the Sindbad Hotel Complex and Conference Center have also been made.
There are also plans to build a giant Ferris wheel akin to the London Eye. Iraq's Tourism Board also is seeking investors to develop a "romantic" island on the River Tigris in Baghdad that was once a popular honeymoon spot for newlywed Iraqis. The project would include a six-star hotel, spa, an 18-hole golf course and a country club. In addition, the go-ahead has been given to build numerous architecturally unique skyscrapers along the Tigris that would develop the city's financial centre in Kadhehemiah.
In October, 2008, the Baghdad Metro resumed service. It connects the center to the southern neighborhood of Dora. In May 2010, a new residential and commercial project nicknamed Baghdad Gate was announced. This project not only addresses the urgent need for new residential units in Baghdad but also acts as a real symbol of progress in the war torn city, as Baghdad has not seen projects of this scale for decades.
The Mustansiriya Madrasah was established in 1227 by the Abbasid Caliph al-Mustansir. The name was changed to Al-Mustansiriya University in 1963. The University of Baghdad is the largest university in Iraq and the second largest in the Arab world.
- Al-Nahrain University
- Al-Mustansiriya University
- Iraqi University
- University of Baghdad
- University of Technology, Iraq
Baghdad has always played an important role in Arab cultural life and has been the home of noted writers, musicians and visual artists. Famous Arab poets and singers such as Nizar Qabbani, Umm Kulthum, Fairuz, Salah Al-Hamdani, Ilham al-Madfai and others wrote beautiful poems and sang for Baghdad.
The dialect of Arabic spoken in Baghdad today differs from that of other large urban centres in Iraq, having features more characteristic of nomadic Arabic dialects (Verseegh, The Arabic Language). It is possible that this was caused by the repopulating of the city with rural residents after the multiple sacks of the late Middle Ages.
Some of the important cultural institutions in the city include:
- Iraqi National Orchestra – Rehearsals and performances were briefly interrupted during the Second Gulf War, but have since returned to normal.
- National Theatre of Iraq – The theatre was looted during the 2003 Invasion of Iraq, but efforts are underway to restore the theatre.
The live theatre scene received a boost during the 1990s when UN sanctions limited the import of foreign films. As many as 30 movie theatres were reported to have been converted to live stages, producing a wide range of comedies and dramatic productions.
Institutions offering cultural education in Baghdad include the Academy of Music, Institute of Fine Arts, and the Music and Ballet school Baghdad. Baghdad is also home to a number of museums which housed artifacts and relics of ancient civilization; many of these were stolen, and the museums looted, during the widespread chaos immediately after United States forces entered the city.
During the 2003 occupation of Iraq, AFN Iraq ("Freedom Radio") broadcast news and entertainment within Baghdad, among other locations. There is also a private radio station called "Dijlah" (named after the Arabic word for the Tigris River) that was created in 2004 as Iraq's first independent talk radio station. Radio Dijlah offices, in the Jamia neighborhood of Baghdad, have been attacked on several occasions.
Baghdad is home to some of the most successful football (soccer) teams in Iraq, the biggest being Al Quwa Al Jawiya (Airforce club), Al Zawra, Al Shurta (Police), and Al Talaba (Students). The largest stadium in Baghdad is Al Shaab Stadium which was opened in 1966. Another, but much larger stadium, is still in the opening stages of construction.
The city has also had a strong tradition of horse racing ever since World War I, known to Baghdadis simply as 'Races'. There are reports of pressures by the Islamists to stop this tradition due to the associated gambling.
Major streets
- Haifa Street
- Salihiya Residential area - situated off Al Sinak bridge in central Baghdad,surrounded by Al- Mansur Hotel in the north and Al-Rasheed hotel in the south
- Hilla Road – Runs from the south into Baghdad via Yarmouk (Baghdad)
- Caliphs Street – site of historical mosques and churches
- Sadoun Street – stretching from Liberation Square to Masbah
- Mohammed Al-Qassim highway near Adhamiyah
- Abu Nuwas Street – runs along the Tigris from the from Jumhouriya Bridge to 14 July Suspended Bridge
- Damascus Street – goes from Damascus Square to the International Airport Road
- Mutanabbi Street – A street with numerous books, named after the 10th century Iraqi poet Al-Mutanabbi
- Rabia Street
- Arbataash Tamuz (14th July) Street (Mosul Road)
- Muthana al-Shaibani Street
- Bor Saeed (Port Said) Street
- Thawra Street
- Al Qanat Street – runs through Baghdad north-south
- Al Khat al Sare'a – Mohammed al Qasim (high speed lane) – runs through Bagdhad, north-south
- Al Sinaa Street (Industry Street) runs by the University of Technology – centre of computers trade in Baghdad
- Al Nidhal Street
- Al Rasheed Street – city centre Baghdad
- Al Jamhuriah Street – city centre Baghdad
- Falastin (Palestine) Street
- Tariq el Muaskar – (Al Rasheed Camp Road)
- Baghdad Airport Road
Sister cities
See also
- Round city of Baghdad
- List of places in Iraq
- Firdos Square - is a public open space in Baghdad and the location of two of the best-known hotels, the Palestine Hotel and the Sheraton Ishtar, which are both also the tallest buildings in Baghdad. The square was the site of the statue of Saddam Hussein that was pulled down by U.S. coalition forces in a widely-televised event during the 2003 invasion of Iraq.
- Operation Imposing Law Baghdad Security Plan
- 1950-1951 Baghdad bombings
- Estimates of total population differ substantially. The Encyclopædia Britannica gives a 2001 population of 4,950,000, the 2006 Lancet Report states a population of 7,216,050 in 2011.
- "Cities and urban areas in Iraq with population over 100,000", Mongabay.com
- "List of largest cities throughout history". Wikipedia.
- Inocencio, Ramy (4 December 2012). "What city has world's best quality of life?". CNN.
- http://www.economist.com/news/middle-east-and-africa/21569077-one-africas-most-miserable-countries-looks-unstable-ever-brink The Central African Republic: On the brink
- "ما معنى اسم مدينة بغداد ومن سماه ؟". Seenjeem.maktoob.com. Retrieved 2010-04-27.
- "ما معنى (بغداد)؟ - تمت الإجابة عنه - Google إجابات". Egabat.google.com. Retrieved 2010-04-27.
- Guy Le Strange, "Baghdad During the Abbasid Caliphate from Contemporary Arabic and Persian", pg 10
- Corzine, Phyllis (2005). The Islamic Empire. Thomson Gale. pp. 68–69.
- Times History of the World. London: Times Books. 2000.
- Wiet, Gastron (1971). Baghdad: Metropolis of the Abbasid Caliphate. Univ. of Oklahoma Press.
- Wiet, pg. 13
- Corzine, Phyllis (2005). The Islamic Empire. Thomson Gale. p. 69.
- Wiet, pg. 12
- "Yakut: Baghdad under the Abbasids, c. 1000CE"
- Wiet, pg. 15
- Hill, Donald R. (1994). Islamic Science and Engineering. Edinburgh: Edinburgh Univ. Press. p. 10. ISBN 0-7486-0457-X.
- Atlas of the Medieval World pg. 78
- "Largest Cities Through History". Geography.about.com. 2009-11-02. Retrieved 2010-04-27.
- Matt T. Rosenberg, Largest Cities Through History.
- George Modelski, World Cities: –3000 to 2000, Washington, D.C.: FAROS 2000, 2003. ISBN 978-0-9676230-1-6. See also Evolutionary World Politics Homepage.
- Trudy Ring, Robert M. Salkin, K. A. Berney, Paul E. Schellinger (1996). International dictionary of historic places, Volume 4: Middle East and Africa. Taylor and Francis. p. 116
- Atlas of the Medieval World pg. 170
- Virani, Shafique N. The Ismailis in the Middle Ages: A History of Survival, A Search for Salvation, (New York: Oxford University Press, 2007), 6.
- Daftary, Farhad. The Isma'ilis: Their History and Doctrines Cambridge: Cambridge University Press, 1990), 205-206.
- Daftary, Farhad. The Isma'ilis: Their History and Doctrines Cambridge: Cambridge University Press, 1990), 206.
- Central Asian world cities, George Modelski
- Ian Frazier, Annals of history: Invaders: Destroying Baghdad, The New Yorker 25 April 2005. p.5
- New Book Looks at Old-Style Central Asian Despotism, EurasiaNet Civil Society, Elizabeth Kiem, April 28, 2006
- "The Fertile Crescent, 1800-1914: a documentary economic history". Charles Philip Issawi (1988). Oxford University Press US. p.99. ISBN 0-19-504951-9
- Suraiya Faroqhi, Halil İnalcık, Donald Quataert (1997). "An economic and social history of the Ottoman Empire". Cambridge University Press. p.651. ISBN 0-521-57455-2
- Cetinsaya, Gokhan. Ottoman Administration of Iraq, 1890-1908. London and New York: Routledge, 2006.
- "The Choice, featuring Lawrence Anthony". BBC radio 4. 2007-09-04. Retrieved 2007-09-04.
- Anthony, Lawrence; Spence Grayham (2007-06-03). Babylon's Ark; The Incredible Wartime Rescue of the Baghdad Zoo. Thomas Dunne Books. ISBN 0-312-35832-6.
- Monthly rainfall for Baghdad (WMO #40650)
- Annual Rainfall Statistics for Baghdad (WMO #40650)
- (AFP) – Jan 11, 2008 (2008-01-11). "Afp.google.com, First snow for 100 years falls on Baghdad". Afp.google.com. Retrieved 2010-04-27.
- "Baghdad Climate Guide to the Average Weather & Temperatures, with Graphs Elucidating Sunshine and Rainfall Data & Information about Wind Speeds & Humidity:". Climate & Temperature. Retrieved 2011-12-25.
- "World Weather Information Service - Baghdad". Retrieved 2010-03-29.
- USA Today http://images.usatoday.com/news/graphics/troop_surge/flash.swf
|url=missing title (help).
- "DefenseLink News Article: Soldier Helps to Form Democracy in Baghdad". Defenselink.mil. Archived from the original on 29 May 2010. Retrieved 2010-04-27.
- "Zafaraniya Residents Get Water Project Update - DefendAmerica News Article". Defendamerica.mil. Retrieved 2010-04-27.
- Frank, Thomas (2006-03-26). "Basics of democracy in Iraq include frustration". USA Today. Retrieved 2010-04-26.
- "DefendAmerica News - Article". Defendamerica.mil. Retrieved 2010-04-27.
- "Democracy from scratch". csmonitor.com. 2003-12-05. Archived from the original on 3 April 2010. Retrieved 2010-04-27.
- "Leaders Highlight Successes of Baghdad Operation - DefendAmerica News Article". Defendamerica.mil. Retrieved 2010-04-27.
- NBC 6 News - 1st Cav Headlines[dead link]
- "Iraqi Airways." Arab Air Carriers Organization. Retrieved on October 19, 2009.
- "Contact Us." Al-Naser Airlines. Retrieved on 13 February 2011. "Main Branch: Al-Karrada , Babil Region - Distrlct 929 [sic] - St21 - Home 46 - Beside Al Jadirya Private Hospital. [...] Iraq- Baghdad."
- ARCADD[dead link]
- 3:48 p.m. ET (2008-08-27). "Baghdad plans to build giant Ferris wheel". MSNBC. Retrieved 2010-04-27.
- "Baghdad Gate". Baghdad Gate. Iraqi News. Retrieved 24 May 2010.
- "Baghdad Investment: Creating (1824) housing units in Baghdad.". Baghdad Governorate Website. 2010. Retrieved 2010-07-09.
- Five women confront a new Iraq | csmonitor.com[dead link]
- "In Baghdad, Art Thrives As War Hovers". Commondreams.org. 2003-01-02. Retrieved 2010-04-27.
- "Gunmen storm independent radio station in latest attack against media in Iraq". International Herald Tribune. 2009-03-29. Retrieved 2010-04-27.
- "Twinning the Cities". City of Beirut. Archived from the original on 2008-02-21. Retrieved 2008-01-13.
Further reading
- By Desert Ways to Baghdad, by Louisa Jebb (Mrs. Roland Wilkins), 1908 (1909 ed) (a searchable facsimile at the University of Georgia Libraries; DjVu & PDF (11.3 MB) format)
- A Dweller in Mesopotamia, being the adventures of an official artist in the Garden of Eden, by Donald Maxwell, 1921 (a searchable facsimile at the University of Georgia Libraries; DjVu & PDF (7.53 MB) format)
- Pieri, Caecilia (2011). Baghdad Arts Deco: Architectural Brickwork, 1920-1950 (1st ed.). The American University in Cairo Press. p. 160. ISBN 978-9774163562.
- "Travels in Asia and Africa 1325-135" by Ibn Battuta.
- "Gertrude Bell: the Arabian diaries,1913-1914." by Bell Gertrude Lowthian, and O'Brien, Rosemary.
- "Historic cities of the Islamic world."by Bosworth, Clifford Edmund.
- "Ottoman administration of Iraq, 1890-1908." by Cetinsaya, Gokhan.
- "Naked in Baghdad." by Garrels, Anne, and Lawrence, Vint.
- "A memoir of Major-General Sir Henry Creswicke Rawlinson." by Rawlinson, George.
|Wikimedia Commons has media related to: Baghdad|
|Wikivoyage has travel information related to: Baghdad|
|Look up Baghdad in Wiktionary, the free dictionary.|
|Wikisource has the text of the 1911 Encyclopædia Britannica article Bagdad (city).|
- Map of Baghdad
- Iraq Image - Baghdad Satellite Observation
- National Commission for Investment in Iraq
- Interactive map
- Iraq - Urban Society
- Envisioning Reconstruction In Iraq
- Description of the original layout of Baghdad
- Ethnic and sectarian map of Baghdad - Healingiraq
- UAE Investors Keen On Taking Part In Baghdad Renaissance Project
- Man With A Plan: Hisham Ashkouri
- Behind Baghdad's 9/11
- Iraq Inter-Agency Information & Analysis Unit Reports, maps and assessments of Iraq from the UN Inter-Agency Information & Analysis Unit |
Samuel Hearne—"The Mungo Park of
Canada"—Perouse complains —The North-West Passage—Indian
guides—Two failures—Third journey successful—Smokes the
calumet—Discovers Arctic Ocean—Cruelty to the Eskimos—Error in
latitude—Remarkable Indian woman—Capture of Prince of Wales
Fort—Criticism by Umfreville.
Such an agitation as that so skilfully
planned and shrewdly carried on by Arthur Dobbs, Esq., could not
but affect the action of the Hudson's Bay Company. The most
serious charge brought against the Company was that, while having
a monopoly of the trade on Hudson Bay, it had taken no steps to
penetrate the country and develop its resources. It is of course
evident that the Company itself could have no reason for refusing
to open up trade with the interior, for by this means it would be
expanding its operations and increasing its profits. The real
reason for its not doing so seems to have been the inertia, not to
say fear, of Hudson's Bay Company agents on the Bay who failed to
mingle with the bands of Indians in the interior.
Now the man was found who was to be equal to
the occasion. This was Samuel Hearne. Except occasional reference
to him in the minutes of the Company and works of the period, we
know little of Samuel Hearne. He was one of the class of men to
which belonged Norton, Kelsey, and others—men who had grown up in
the service of the Company on the Bay, and had become, in the
course of years, accustomed to the climate, condition of life, and
haunts of the Indians, thus being fitted for active work for the
Hearne became so celebrated in his inland expeditions, that the
credit of the Hudson's Bay Company leaving the coast and venturing
into the interior has always been attached to his name. So
greatly, especially in the English mind, have his explorations
bulked, that the author of a book of travels in Canada about the
beginning of this century called him the "Mungo Park of Canada."
In his "Journey," we have an account of his earlier voyages to the
interior in search of the Coppermine River. This book has a
somewhat notable history. In the four-volume work of La Perouse,
the French navigator, it is stated that when he took Prince of
Wales Fort on the Churchill River in 1782, Hearne, as governor of
the fort, surrendered it to him, and that the manuscript of his
"Journey" was seized by the French commander. It was returned to
Hearne on condition that it should be published, but the
publication did not take place until thirteen years afterwards. It
is somewhat amusing to read in Perouse's preface (1791) the
complaint that Hearne had not kept faith with him in regard to
publishing the journal, and the hope is expressed that this public
statement in reminding him of his promise would have the desired
effect of the journal being published.
Four years afterwards Hearne's "Journey"
appeared. A reference to this fine quarto work, which is well
illustrated, brings us back in the introduction to all the
controversies embodied in the work of Dobbs, Ellis, Robson, and
the "American Traveller."
Hearne's orders were received from the
Hudson's Bay Company, in 1769, to go on a land expedition to the
interior of the continent, from the mouth of the Churchill as far
as 70 deg. N. lat., to smoke the calumet of peace with the
Indians, to take accurate astronomical observations, to go with
guides to the Athabasca country, and thence northward to a river
abounding with copper ore and "animals of the fur kind," &c.
It is very noticeable, also, that his
instructions distinctly tell him" to clear up the point, if
possible, in order to prevent further doubt from arising hereafter
respecting a passage out of Hudson Bay into the Western Ocean, as
hath lately been represented by the 'American Traveller.'" The
instructions made it plain that it was the agitation still
continuing from the days of Dobbs which led to the sending of
Hearne to the north country.
Hearne's first expedition was made during the
last months of the year 1769. It is peculiarly instructive in the
fact that it failed to accomplish anything, as it gives us a
glimpse of the difficulties which no doubt so long prevented the
movement to the interior. In the first place, the bitterly severe
months of November and December were badly chosen for the time of
the expedition. On the sixth day of the former of these months
Hearne left Prince of Wales Fort, taking leave of the Governor,
and being sent off with a salute of seven guns. His guide was an
Indian chief, Chawchinahaw. Hearne ascertained very soon, what
others have found among the Indians, that his guide was not to be
trusted; he "often painted the difficulties in the worst colours"
and took every method to dishearten the explorer. Three weeks
after starting, a number of the Indians deserted Hearne.
Shortly after this mishap, Chawchinahaw and
his company ruthlessly deserted the expedition, and two hundred
miles from the fort set out on another route, "making the woods
ring with their laughter." Meeting other Indians, Hearne purchased
venison, but was cheated, while his Indian guide was feasted. The
explorer remarks:—"A sufficient proof of the singular advantage
which a native of this country has over an Englishman, when at
such a distance from the Company's factories as to depend entirely
on them for subsistence."
Hearne arrived at the fort after an absence
of thirty-seven days, as he says, "to my own mortification and the
no small surprise of the Governor." Hearne was simply illustrating
what has been shown a hundred times since, in all foreign regions,
viz., native peoples are quick to see the inexperience of men raw
to the country, and will heartlessly maltreat and deceive them.
However, British officers and men in all parts of the world become
at length accustomed to dealing with savage peoples, and after
some experience, none have ever equalled British agents and
explorers in the management and direction of such peoples.
Early in the following year Hearne plucked up
courage for another expedition. On this occasion ho determined to
take no Europeans, but to trust to Indians alone. On February
23rd, accompanied by five Indians, Hearne started on his second
journey. Following the advice of the Governor, the party took no
Indian women with them, though Hearne states that this was a
mistake, as they were "needed for hauling the baggage as well as
for dressing skins for clothing, pitching our tent, getting
firing, &c." During the first part of the journey deer were
plentiful, and the fish obtained by cutting holes in the ice of
the lakes were excellent.
Hearne spent the time of the necessary delays
caused by the obtaining of fish and game in taking observations,
keeping his journal and chart, and doing his share of trapping.
Meeting, as soon as the spring opened, bands of Indians going on
various errands, the explorer started overland. He carried sixty
pounds of burden, consisting of quadrant, books and papers,
compass, wearing apparel, weapons and presents for the natives.
The traveller often made twenty miles a day over the rugged
Meeting a chief of the Northern Indians going
in July to Prince of Wales Fort, Hearne sent by him for ammunition
and supplies. A canoe being now necessary, Hearne purchased this
of the Indians. It was obtained by the exchange of a single knife,
the full value of which did not exceed a penny. In the middle of
this month the party saw bands of musk oxen. A number of these
were killed and their flesh made into pemmican for future use.
Finding it impossible to reach the Coppermine during the season,
Hearne determined to live with the Indians for the winter.
The explorer was a good deal disturbed by
having to give presents to Indians who met him. Some of them
wanted guns, all wanted ammunition, iron-work, and tobacco; many
were solicitous for medicine; and others pressed for different
articles of clothing. He thought the Indians very inconsiderate in
On August 11th the explorer had the
misfortune to lose his quadrant by its being blown open and broken
by the wind. Shortly after this disaster, Hearne was plundered by
a number of Indians who joined him.
He determined to return to the fort.
Suffering from the want of food and clothing, Hearne was overtaken
by a famous chief, Matonabbee, who was going eastward to Prince of
Wales Fort. The chief had lived several years at the fort, and was
one who knew the Coppermine. Matonabbee discussed the reasons of
Hearne's failure in his two expeditions. The forest philosopher
gave as the reason of these failures the misconduct of the guides
and the failure to take any women on the journey. After
maintaining that women were made for labour, and speaking of their
assistance, said Matonabbee, "women, though they do everything,
are maintained at a trifling expense, for as they always stand
cook, the very licking of their fingers in scarce times is
sufficient for their subsistence." Plainly, the northern chief had
need of the ameliorating influence of modern reformers. In company
with the chief, Hearne returned to the fort, reaching it after an
absence of eight months and twenty-two days, having, as he says,
had "a fruitless or at least an unsuccessful Journey."
Hearne, though beaten twice, was determined
to try a third time and win. He recommended the employment of
Matonabbee as a guide of intelligence and experience. Governor
Norton wished to send some of the coast Indians with Hearne, but
the latter refused them, and incurred the ill-will of the
Governor. Hearne's instructions on this third Journey were "in
quest of a North-West Passage, copper-mines, or any other thing
that may be serviceable to the British nation in general, or the
Hudson's Bay Company in particular." The explorer was now
furnished with an Elton's quadrant.
This third Journey was begun on December 7th,
1770. Travelling sometimes for three or four days without food,
they were annoyed, when supplies were secured, by the chief
Matonabbee taking so ill from over-eating that he had to be drawn
upon a sledge. Without more than the usual incidents of Indian
travelling, the party pushed on till a point some 19 deg. west of
Churchill was reached, according to the calculations of the
explorer. It is to be noted, however, that Hearne's observations,
measurements, and maps, do not seem to be at all accurate.
Turning northward, as far as can be now made
out, about the spot whore the North-West traders first appeared on
their way to the Churchill River, Hearne went north to his
destination.1 His Indian guides now formed a large war party from
the resident Indians, to meet the Eskimos of the river to which
they were going and to conquer them.
The explorer announces that having left
behind "all the women, children, dogs, heavy baggage, and other
encumbrances," on June 1st, 1771, they pursued their journey
northward with great speed. On June 21st the sun did not set at
all, which Hearne took to be proof that they had reached the
Arctic Circle. Next day they met the Copper Indians, who welcomed
them on hearing the object of their visit.
Hearne, according to orders, smoked the
calumet of peace with the Copper Indians. These Indians had never
before seen a white man. Hearne was considered a great curiosity.
Pushing on upon their long journey, the explorers reached the
Coppermine River on July 13th. Hearne was the witness of a cruel
massacre of the Eskimos by his Indian allies, and the seizure of
their copper utensils and other provisions, and expresses disgust
at the enormity of the affair. The mouth of the river, which flows
into the Arctic Ocean, was soon reached on July 18th, and the tide
found to rise about fourteen feet.
Hearne seems in the narrative rather
uncertain about the latitude of the mouth of the Coppermine River,
but states that after some consultation with the Indians, he
erected a mark, and took possession of the coast on behalf of the
Hudson's Bay Company.
In Hearne's map, dated July, 1771, and
purporting to be a plan of the Coppermine, the mouth of the river
is about 71 deg. 54' N. This was a great mistake, as the mouth of
the river is somewhere near 68 deg. N. So great a mistake was
certainly unpardonable. Hearne's apology was that after the
breaking of his quadrant on the second expedition, the instrument
which he used was an old Elton's quadrant, which had been knocking
about the Prince of Wales Fort for nearly thirty years.
Having examined the resources of the river
and heard of the mines from which the Copper Indians obtained all
the metal for the manufacture of hatchets, chisels, knives, &c,
Hearne started southward on his return journey on July 18th.
Instead of coming by the direct route, he went with the Indians of
his party to the north side of Lake Athabasca on December 24th.
Having crossed the lake, as illustrating the loneliness of the
region, the party found a woman who had escaped from an Indian
band which had taken her prisoner, and who had not seen a human
face for seven months, and had lived by snaring partridges,
rabbits, and squirrels. Her skill in maintaining herself in lonely
wilds was truly wonderful. She became the wife of one of the
Indians of Hearne's party. In the middle of March, 1772, Hearne
was delivered a letter, brought to him from Prince of Wales Fort
and dated in the preceding June. Pushing eastward, after a number
of adventures, Hearne reached Prince of Wales Fort on June 30th,
1772, having been absent on his third voyage eighteen months and
twenty-three days. Hearne rejoices that he had at length put an
end to the disputes concerning a North-West Passage through Hudson
Bay. The fact, however, that during the nineteenth century this
became again a living question shows that in this he was mistaken.
The perseverance and pluck of Hearne have
impressed all those who have read his narrative. He was plainly
one of the men possessing the subtle power of impressing the
Indian mind. His disasters would have deterred many men from
following up so difficult and extensive a route. To him the
Hudson's Bay Company owes a debt of gratitude. That debt consists
not in the discovery of the Coppermine, but in the attitude
presented to the Northern Indians from the Bay all the way to Lake
Athabasca, Hearne does not mention the Montreal fur traders, who,
in the very year of his return, reached the Saskatchewan and were
stationed at the Churchill River down which ho passed.
First of white men to reach Athapuscow, now
thought to have been Great Slave Lake, Samuel Hearne claimed for
his Company priority of trade, and answered the calumnies that his
Company was lacking in energy and enterprise. Ho took what may be
called "seizen" of the soil for the English traders. We shall
speak again of his part in leading the movement inland to oppose
the Nor'-Westers in the interior. His services to the Hudson's Bay
Company received recognition in his promotion, three years after
his return home from his third voyage, to the governorship of the
Prince of Wales Fort. To Hearne has been largely given the credit
of the new and adventurous policy of the Hudson's Bay Company, -
Hearne does not, however, disappear from public notice on his
promotion to the command of Prince of Wales Fort. When war broke
out a few years later between England and France, the latter
country, remembering her old successes under D'lber-ville on
Hudson Bay, sent a naval expedition to attack the forts on the
Bay. Umfreville gives an account of the attack on Prince of Wales
Fort on August 8th and 9th, 1772. Admiral de la Perouse was in
command of these war vessels, his flagship being Le Sceptre, of
seventy-four guns. The garrison was thought to be well provided
for a siege, and La Perouse evidently expected to have a severe
contest. However, as he approached the fort, there seemed to be no
preparations made for defence, and, on the summons to surrender,
the gates were immediately thrown open.
Umfreville, who was in the garrison and was
taken prisoner on this occasion, speaks of the conduct of the
Governor as being very reprehensible, but severely criticizes the
Company for its neglect. He says:—"The strength of the fort itself
was such as would have resisted the attack of a more considerable
force; it was built of the strongest materials, the walls were of
great thickness and very durable (it was planned by the ingenious
Mr. Robson, who went out in 1742 for that purpose), it having been
forty years in building and attended with great expense to the
Company. In short, it was the opinion of every intelligent person
that it might have made an obstinate resistance when attacked, had
it been as well provided in other respects; but through the
impolitic conduct of the Company, every courageous exertion of
their servants must have been considered as imprudent temerity;
for this place, which would have required four hundred men for its
defence, the Company, in its consummate wisdom, had garrisoned
with only thirty-nine."
In this matter, Umfreville very plainly shows
his animus to the Company, but incidentally he exonerates Hearne
from the charge of cowardice, inasmuch as it would have been
madness to make defence against so large a body of men. As has
been before pointed out, we can hardly charge with cowardice the
man who had shown his courage and determination in the three
toilsome and dangerous journeys spoken of; rather would we see in
this a proof of his wisdom under unfortunate circumstances. The
surrender of York Factory to La Perouse twelve days afterwards,
without resistance, was an event of an equally discouraging kind.
The Company suffered great loss by the surrender of these forts,
which had been unmolested since the Treaty of Utrecht. |
"Anabaptist" is actually a Greek word meaning "rebaptizer," used in church Latin from the 4th century onward, and appearing at least as early as 1532 in the English, seldom used in 16th-century German or Dutch, where the translation Wiedertäufer and Wederdooper is used from the beginning of Anabaptist history in 1525. It was never used by the Anabaptists themselves but often vigorously objected to by them because of the opprobrium and criminal character attached to the name. Its introduction and constant use by the enemies of the Anabaptists can best be explained by the fact that the imperial law code from Justinian's time (A.D. 529) on, made rebaptism one of the two heresies penalized by death, the other being Antitrinitarianism. Thus to classify the Reformation radicals as "Anabaptists" made them at once legally subject to condemnation and execution, although it still remained necessary for each local jurisdiction to implement the basic code. (Thus Zürich did not decree the death penalty for Anabaptists until 1526.) The first imperial mandate (4 January 1528, Speyer) against the Anabaptists specifically grounds the required suppression on the ancient imperial law as follows: "Since in both ecclesiastical and civil law Anabaptism (der Widertauf) is forbidden under severe penalties, and since the imperial code decrees and orders, on pain of the highest penalty of death, that no one shall have himself baptized a second time or baptize another . . . . "
The very first literary attacks on Anabaptism (Zwingli's Von der Taufe, von der Wiedertaufe, und von der Kindertaufe of May 1525, and Oecolampadius' Ein Gespräch etlicher Predicanten zu Basel gehalten mit etlichen Bekennern de, Wiedertaufe, also of 1525) use the term Wiedertäufer. The Latin writings likewise use Anabaptists (e.g., Faber's Adversus Doctorem Balthasarum Pacimontanium, Anabaptistarum Nostri Saeculi, Primum Authorem, Orthodoxae Fidei Catholica Defensio of 1528). However, in Zwingli's testimony against the local Anabaptists before the Zürich court in March 1525 he uses the term Täufer exclusively. Strangely enough, Zwingli uses the term "Catabaptist" in his major attack against the group in 1527, In Catabaptistarum Strophas Elenchus, instead of "Anabaptist."
It has sometimes been assumed that the evil connotation of the epithet Anabaptist is associated primarily with the dreadful Münster episode of 1534-1535. However, the fairly extensive polemic literature of the period before that time, written by Luther, Melanchthon, Zwingli, Bader, Rhegius, Faber, Bugenhagen, Menius, Bullinger, and others, gives abundant evidence that it was a designation of severest reproach and condemnation long before Münster. The Augsburg Confession of 1530 condemns the "Anabaptists" specifically in three articles, though in part based on misinformation. Abundant citations could be given showing that the term "Anabaptist" in all its forms and translations was always essentially one of condemnation as of grievous heresy and crime. More, in his Confutation of Tindale's Works (1532), speaks of "pernitious and Anabaptistical opinions." This completely evil connotation of the name, which makes it truly an opprobrious epithet, carried through the 16th century and on down through the following centuries until modern times. It is this sense of condemnation and execration which has made some modern historians, particularly Mennonites, hesitate. to use it, but usage is gradually overcoming the objectionable sense.
An illustration of the strong objection by those, dubbed Anabaptists to its application to them is the title of Dirk Philips' Dutch tract written in 1545 but first published as the fourth part of his Enchiridion in 1564, which reads as follows: Een Apologia, ofte verantwoordinghe, dat wy (die van de werelt met grooten onrecht Anabaptisten gescholden worden) gheen wederdoopers noch sectemakers en zijn; maer dat wy een zijn met de rechte Ghemeynte Gods die van aenbeginne gheweest is. (An Apology or Reply, that we who are by the world with great injustice accusingly called Anabaptists are no re-baptizers nor sect-makers, but that we are one with the true church of God which has been from the beginning.)
The original use of "Anabaptist" in the 4th and following centuries was to refer to the rebaptism of those who had been baptized by heretics, or of those who had been baptized by bishops who had temporarily and partially recanted under persecution. There was considerable controversy over both points in North Africa; over the former from Tertullian's time on (A.D. 200) and over the latter in Augustine's time (Donatist controversy). In both cases the Roman bishop's position won out, namely, that rebaptism should not be required nor permitted. Those who insisted on rebaptism were in effect repudiating therefore the authority of the church.
The Anabaptists of the Reformation period, however, did not repudiate infant baptism because they denied the validity of office of the bishop or the authority of the church (although they did in fact deny both) but rather because they denied the readiness of an infant to receive baptism on New Testament terms. They called for baptism only on confession of faith and commitment to discipleship by the candidate. They denied that infant baptism was baptism at all and hence denied that they were "rebaptizers." However, their real objection to the name "'Anabaptist" was not this minor technical one; it was rather their refusal to be classed as heretics and to be reckoned as not being the true church. Their intensity of feeling on this must be understood in the light of their deep conviction that they were the true church and that the Roman Catholic and Protestant churches were the false churches. Naturally also they did not wish to be classed as heretics subject to the death penalty merely on the basis of an epithetical identification with the Anabaptists of earlier centuries whom the imperial law condemned to death. They wished to stand on their own faith and to have their testimony and doctrine received on its own merits.
Although the meaning put into "Anabaptist" before 1535 was bad enough, additional overtones were added after Münster. Already frightened by the rapid spread as well as obstinate steadfastness and evident spiritual power of the movement, and already firmly believing the Anabaptists to be a threat to the existing order and social stability, the leaders of church and state were now sure of it for the Münsterites were actually militant revolutionaries and perverters of Christian ethics.
Hitherto the Anabaptists had lived irreproachable lives, but now the scandalous behavior of the King of Münster and his henchmen was known to all. So the invective against the Anabaptists now rises to a shrill crescendo. The Protestants in particular were concerned to vindicate themselves of the Catholic charge of complicity in and responsibility for the Anabaptists by going to extreme lengths of condemnation. The epithet "Anabaptist" was thus filled with even more venom than before, if that could be possible. It became the synonym for everything dangerous to church and state, much like "Bolshevik" or "Communist" in 1950s America.
Furthermore, the epithet was used indiscriminately of all types of left-wing groups, whereby the sins of the worst were applied to all. In retrospect Thomas Müntzer, the leader of the Peasants' Revolt in 1525, was now dubbed an Anabaptist and his sins were added to the others. In fact, he came to be thought of as the originator and most typical leader of the movement, even though he never practiced nor taught rebaptism, and had no connection with the true Anabaptist movement.
The Anabaptists themselves used no common name, indeed they were not a unified organized movement throughout, although the Swiss-South German, Dutch-North German, and Hutterite wings were soon separately organized and disciplined. Their most common self-designation was "Brethren." Because of the strong leadership of Jakob Hutter (d. 1535) among the Moravian Anabaptists, who adopted community of goods, this group was soon called "Hutterisch" or the "Hutterian Brethren," while the non-communist group being originally of Swiss origin was called "Swiss Brethren," even though they lived in many places outside of Switzerland such as the lower Rhine region. In Holland after 1545 the group came to be called "Mennists" after their chief leader Menno Simons, a name which gradually developed into "Mennonist" and then "Mennonit," although early in the 17th century "Doopsgezind" (German, "Taufgesinnt") came into use in Holland and ultimately superseded "Mennonit." Thus at least one syllable of the original "Wederdooper" is retained in the modern designation of the descendants of the original Dutch Anabaptists.
The 17th-century English Baptists took the major part of "Anabaptists" for their name and passed it on down to their 10,000,000 modern spiritual descendants. The 19th-century Baptists of Germany however did not take "Täufer," the proper translation, for their name, but instead, "Baptisten." It remained for modern German scholars to adopt the term "Täufer" as the designation for the Anabaptists of the Reformation period, thus taking up Zwingli's early 1525 word instead of 'Wiedertäufer." -- HSB
The German word Täufer has a definite meaning and is applied exclusively today to those evangelical Anabaptists who represent the ancestors of today's Mennonites all over the world. German church historiography has generally abandoned the terms Anabaptisten and Wiedertäufer; Dutch historians use the terms Doopers or Doopsgezinden for the German Täufer. In English, American, and French church historiography, however, the term Anabaptist (Dutch Wederdooper, German Wiedertäufer, French Anabaptiste) is used in a much broader and more inclusive sense to cover all types of radicals of the Reformation period. Basically the word "Anabaptist" indicates nothing more than the rejection of infant baptism (on whatever ground) and the consequent practice of adult baptism or baptism on confession of faith. This general principle was, in the 16th century, held by many different groups which can scarcely be thrown together indiscriminately under the heading "Anabaptist" in the sense that it is used in this encyclopedia (although from the 16th to the 18th centuries it was almost universally so used in Europe). For instance, the Socinian Polish church accepted this principle without belonging to what generally is called Anabaptism. The same is true of scattered groups in England in the 16th century, which sometimes are called Anabaptists, or with scattered individuals of the "enthusiastic" fringe who have only a very loose connection with the main body of evangelical Anabaptists, or the fanatical revolutionary Münsterites of 1534-1535. Professor Bainton of Yale, for instance, almost identifies the term "Anabaptist" with what he calls Left-Wing Protestantism in general, that radical wing of widely varied types which pressed for complete separation of state and church, and at the same time for a believers' or gathered church. (In German church history, mainly through the work of E. Troeltsch, The Social Teachings of the Christian Churches, 1931, first German edition, 1912, the term "sect" is used for these groups in general. Unfortunately it is not applicable in English and American church history.) It is quite obvious that this interpretation of the concept of Anabaptism is much too broad to be usefully applied to any one distinct body. Not all Protestant groups which practice adult baptism can be classified as Anabaptists, as can be clearly seen, for instance, in the distinction between Anabaptists and Baptists. But even so, the term still needs a more precise definition whereby a distinction must also be drawn between the stillen Täufer or evangelical Anabaptists (see Church History, 1940, 362) who accept the principle of nonresistance, and all other related groups which decline nonresistance, and accept the "sword" as a positive instrument. To this latter group may be counted the millennialist Münsterites and their partisans in Holland and elsewhere. (Troeltsch suggested the term "Taborites" for this type. In English church history the "Fifth Monarchy Men" would fall under this category.)
As a matter of fact, the phenomenon of Anabaptism has been subjected to manifold and widely varying interpretations, and even in the 1950s, with greatly improved factual knowledge, there was no complete agreement among church historians regarding its understanding and definition.
Anabaptism certainly does not simply mean (a) the refusal of infant baptism for whatever reason. Even the reformers themselves, Luther and Zwingli, admitted in the earlier years of their work that infant baptism is without Scriptural basis. (b) Nor is Anabaptism the same as fanaticism (Schwärmertum, Schwarmgeisterei). This is the traditional confusion mainly among Lutherans (e.g., Karl Holl) because Luther himself called all his opponents Schwärmer, i.e., people who have replaced Biblical theology with personal inspiration, special illumination, or private revelations. It is quite obvious that most Anabaptists were far from any such fanaticism (which they hated just as much as did Luther). Their strict Biblicism is beyond doubt; it was at once their strength and their limitation. (c) Anabaptism has hardly anything in common with traditional millennialism, even though certain millennialists favored and still favor adult baptism (e.g., the Münsterites, the Fifth Monarchy Men, later some Pietists). Millennialism was for Anabaptists always but a marginal idea. Finally (d) Anabaptism should not be confused with Antitrinitarianism as has so often been done and still occasionally is done (Dunn-Borkowski and Wilbur: see Antitrinitarianism). The Polish Socinians, though practicing adult baptism, declined any connection with western Anabaptism; and yet, Wilbur still considers many early Anabaptists as forerunners of Unitarianism. Thus all these interpretations expanded the term "Anabaptist" into a general concept of Free Church Protestantism, by which the concept of genuine Anabaptism loses its character and its meaningfulness.
We now turn to modern attempts at positive definition or delineation of the concept of Anabaptism.
It is quite clear that this rapid survey by no means exhausts the problem involved in a delineation of the idea of evangelical Anabaptism. The present analysis has been more negative than positive. It needs, therefore, ample supplementation. H. S. Bender's Anabaptist Vision and his analysis of the church concept of the Brethren (Mennonite Quarterly Review 18, 1944, 67-88, and 19, 1945, 90-100) might well supply it. (R. Friedmann, "Conception of an Anabaptist," Church History 9, 1940, 341-365.) -- RF
After summarizing the perception of what Anabaptism was after the revision carried out by Harold S. Bender and others, Robert Friedmann asserts in volume 1 of the Mennonite Encyclopedia above that Anabaptists were "nearer to the spirit of Christ's exemplary life and teaching. . ." than were Protestants, Schwärmer ("enthusiasts"), and millennialists, and that Bender's "Anabaptist vision" might supply "delineation of the idea of evangelical Anabaptism." This had, in fact, already happened by the time Friedmann wrote. "Evangelical Anabaptism" was identified by insistence on discipleship as the essence of Christianity, on the church as a brotherhood, and on an ethic of love and nonresistance. This became the normative description of Anabaptism. In this view, evangelical Anabaptism arose with the Swiss Brethren, and by transmission became part of Netherlands Anabaptism and of the Hutterites. Thus was Anabaptism given unity and clearly distinguished from Catholicism, from Protestantism, and from other 16th-century dissenting groups.
A revision of this portrayal began around 1960. Heinold Fast warned in a 1967 article that the Mennonite revision of four centuries of negative historiography was too tidy, too ideal, and that reaction would come. Indeed, reaction was already under way as part of a major shift in Reformation studies from systematic theology to history of ideas and from confessional history to social history. This shift strongly modified the traditional confessional (Lutheran, Reformed, Catholic, Mennonite, etc.) orientation and opened the door to consideration of the dialectic between social and political developments, on the one hand, and the development of theological positions in the various reinterpretations of Christian faith, on the other. For Anabaptist studies it meant the entry into the field of a number of non-Mennonite historians who studied Anabaptism not as the ancestral movement of the 20th-century church communion, but as part of the general history of western Europe in the 16th century. It produced a new picture of Anabaptism not only socially but also in theology.
The word Anabaptism is normally used in the late 20th century to denote the mosaic of groupings of dissenters without at the same time making claims to uniformity. In his 1972 work Anabaptists and the Sword, James M. Stayer used the term with great care in order to avoid giving the impression that he was writing about a single unified movement across Europe. He wrote about Anabaptists and defined them as those who rebaptized persons already baptized in infancy. Walter Klaassen had already used this definition in his Oxford dissertation in 1960. Calvin Pater (1984) broadened the definition by including those who, before 1525, rejected the baptism of infants, but this is perhaps too broad to be useful. These definitions were meant to avoid such confessional definitions as "evangelical Anabaptism." Thus, all those rebaptizers who have in the past been classified as Schwärmer (Melchior Hoffman), spiritualists (Hans Denck), and revolutionaries (the Münsterites) are now considered to be genuine Anabaptists.
Despite the variety of viewpoints among 16th-century Anabaptists, despite important differences of nuance even where Anabaptists appear to be similar, one may hazard to identify some themes held in common following the crystallization of the movement between 1527 and 1540. (1) All shared a basically synergistic view of salvation (human and divine "cooperation"). Justification was seen as progression in holiness; the ethic of the Sermon on the Mount was the guide to it. (2) Baptism was considered to be the sign of lay emancipation from clerical control and the spiritual enfranchisement of lay people (priesthood of all believers). (3) Anabaptists developed a Gemeindechristentum centered on the congregation, in contrast to the clericalized territorial churches, both Catholic and Protestant.
However, the open definition of Anabaptist now in use emphatically does not imply uniformity. Anabaptism was pluralistic. Claus-Peter Clasen identified six major groupings often hostile to each other, and then cited contemporary literature to show that there were actually many more (1972). The 1975 article "From Monogenesis to Polygenesis," by James M. Stayer, Werner O. Packull, and Klaus Deppermann, has become the accepted statement on Anabaptist plurality. The essay disputed the older view that Anabaptism had its origins solely in Zürich, and that Swiss Brethren Anabaptism was transmitted to South Germany and Austria and to the Netherlands and North Germany, where it developed into the Hutterian and Mennonite branches respectively. The authors showed that each of the three in fact had a distinctive character and therefore a distinct source. For South German-Austrian Anabaptism it was a diluted form of Rhineland mysticism (Packull, 1977). Social unrest and the apocalyptic visions of Melchior Hoffman put their stamp on Netherlands Anabaptism (Deppermann, 1979; 1987). Swiss Anabaptism arose out of Reformed congregationalism (Stayer, 1975).
Numerous individual studies demonstrating links and relationships between Anabaptists have gradually led to the abandonment of the Schleitheim Confession as a norm for all "true" Anabaptists. As long ago as 1956 Frank J. Wray showed that Pilgram Marpeck had borrowed the bulk of his Vermanung (1542) from the despised Münster theologian Bernhard Rothmann's Bekenntnisse of 1533. Quite as surprising was the demonstration that Melchior Hoffman's commentary on the Apocalypse (1530) was used by the Hutterites soon after Hoffman wrote it, but without acknowledgement of authorship (Packull, 1982).
David Joris had for long been a pariah, especially for North American Mennonite historians, and was condemned by relative inattention. Two major dissertations have shown Joris to have been the most important Anabaptist leader in the Netherlands before 1540, even more so than Menno Simons (Zijlstra, 1983; Waite, 1986). He was an influential figure in Anabaptism's consolidation period following the fall of Münster.
Most significant is the integrative rewriting of the history of Anabaptism in the Netherlands. Melchior Hoffman is acknowledged as its progenitor, as the person who gave the movement its basic apocalyptic stamp. The differences that appear along the "peaceful" to "revolutionary" spectrum can be accounted for by differing nuances in the era's widespread apocalyptic expectation (Klaassen, 1986). A clear line stretches from Hoffman to Rothmann, Menno Simons, and David Joris on apocalyptic anticipation. Very similar formulations of apocalyptic views on secular government and on the incarnation are found in Hoffman, Rothmann, and Menno Simons (Voolstra, 1982; Stayer, 1972, 1978, 1986). The work of these scholars has therefore shown that in some important respects there was a single movement from Hoffman to Menno.
George H. Williams' massive volume, The Radical Reformation (1962), presented for the first time a comprehensive portrayal of radical dissent in the 16th century. Moving across his stage are the whole cast of characters, from Andreas Karlstadt (Carlstadt) and the Zwickau Prophets through Conrad Grebel and the St. Gall fanatics, Hans Hut and Jakob Hutter, Menno Simons and the Batenburgers, to Michael Servetus and Faustus Socinus, with all the intricate linkages between them.
Clear links between Thomas Müntzer and Hans Hut had been established by Grete Mecenseffy (1956) and Walter Klaassen (1960, 1962), but the most important study on this relationship was done by Gottfried Seebaß (1972). His work and the work of Werner Packull (1977) established beyond question the formative influence of Thomas Müntzer on South German Anabaptism, both in its mysticism and its apocalyptic cast. By means of the thesis that mystical theology was a theology of dissent in the 16th century, Steven Ozment linked the Anabaptists Hans Denck and Hans Hut with Müntzer, Sebastian Franck, and others. The central thrust of this mysticism was that ultimately God could communicate his will to men and women directly in disregard of ecclesiastical channels, a kind of democratization of revelation. Later Calvin Pater (1984) convincingly showed that Andreas Karlstadt significantly influenced Swiss Anabaptism. Versions of Karlstadt's view and use of Scripture, his doctrine of the church, and his views on baptism, all found their way into Anabaptism.
Building on earlier work, Hans-Jürgen Goertz (1980) offered the first extended discussion of the thesis that anticlericalism was a prime motive for dissent, and that this factor provided close links between Anabaptism and other movements of the "common man," such as the peasant uprisings of 1524-26, also extensively motivated by anticlericalism. Among Anabaptists this expressed itself in contrasts between the Good Shepherd and the self-indulgent clergy, the simple reading of Scripture and its use as a means of oppression, and the improvement of life and the fruitless life of the new Protestant teachers of justification by faith alone. Other expressions of anticlericalism were the involvement of the Zürich radicals in opposition to tithes and the demand for congregational autonomy. Both central issues for peasants have been clearly documented by Haas (1975) and Stayer (1975A). Werner Packull (1985) and Arnold Snyder (1984, 1985) have provided further evidence for these links.
Finally, the relationship of Anabaptists to Caspar Schwenckfeld has also been extensively studied in recent years. Neal Blough's work on Pilgram Marpeck (1984) demonstrates dependence of Marpeck on Schwenckfeld especially relating to their understanding of the Incarnation. A complex set of relationships of Schwenckfeld with Melchior Hoffman and Pilgram Marpeck was described by R. Emmet McLaughlin (1985). The lure of Schwenckfeld's spiritualism for South German Anabaptists was clearly shown by George H. Williams in his Radical Reformation. Anabaptists and Schwenckfeld agreed on many important issues (Klaassen, 1986).
Anabaptists must now therefore be seen as an integral part of the larger phenomenon of religious and social dissent in 16th-century Europe from the Zwickau Prophets to Sébastian Castellion.
Anabaptism arose out of the religious and social ferment of the Reformation period. That Anabaptists everywhere should have been influenced in numerous ways by the Reformers, is established. That they were the most consistent Protestants carrying the reforms of Luther and Zwingli to their logical conclusion as earlier Mennonite interpreters such as Cornelius H. Wedel and John Horsch held, is a view that can no longer be sustained. For Anabaptists differed with the major Reformers both on the principles of sola scriptura (by Scripture alone), and even more radically on sola fidei (by faith alone). Because of their profound concern for ethics, they adopted variants of a synergistic soteriology which bore resemblance to some late medieval views. But there was also ambivalence among Anabaptists as to whether they were reformist or restitutionist (Wray, 1954; Meihuizen, 1970). As Hans J. Hillerbrand (1971) pointed out, restitutionists have great difficulty dealing with recent history. This point has also been made by Dennis Martin (1987), who argued that restitutionists, in contrast to reformers, can really build no lasting traditions since their revolt against a corrupt immediate past makes them suspicious even of any new institutions or traditions they may establish.
Finally, the question as to whether Anabaptism was medieval or modern has been vigorously debated. The link of Anabaptism to mysticism, its synergistic soteriology, and its version of imitatio Christi, all point to pre-Reformation forms of piety (Ozment, 1972; Davis, 1974; Packull, 1977). Alternatively, it has been argued that Anabaptism was the true harbinger of modernity in its emphasis on voluntarism, toleration, and pluralism in religion (Bender, 1955, Zeman, 1976). The early Swiss Brethren, claimed Fritz Blanke, were a vanguard striving toward a new dawn (Blanke, 1961). A carefully nuanced statement on this subject describes social tendencies in Anabaptism that moved in the direction of modernity (Goertz, 1985).
Bender, Harold S. "The Anabaptists and Religious Liberty in the 16th Century." Mennonite Quarterly Review 29 (1955): 83-100.
Blanke, Fritz Blanke. Brothers in Christ. Scottdale, 1961.
Blough, Neal. Christologie Anabaptiste: Pilgram Marpeck et L'humanité du Christ. Genève: Éditions Labor et Fides, 1984.
Clasen, Claus-Peter. Anabaptism: A Social History 1525-1618. Ithaca: Cornell U. Press, 1972.
Davis, Kenneth R. Anabaptism and Asceticism. Scottdale, 1974.
Deppermann, Klaus. Melchior Hoffman: Soziale Unruhen und apokalyptische Visionen im Zeitalter der Reformation. Göttingen: Vandenhoeck und Ruprecht, 1979; English translation, Edinburgh: T. and T. Clark, 1987.
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Pater, Calvin A. Karlstadt as the Father of the Baptist Movements; The Emergence of Lay Protestantism. Toronto: U. of Toronto Press, 1984.
Seebaß, Gottfried. "Müntzers Erbe: Werk, Leben und Theologie des Hans Hut." unpub. Habilitationsschrift, Erlangen, 1972.
Snyder, C. Arnold. The Life and Thought of Michael Sattler. Scottdale, 1984.
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Stayer, James M. "Die Anfänge des schweizerischen Täufertums im reformierten Kongregationalismus," in Umstrittenes Täufertum (1975): 19-49.
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Stayer, James M. "Was Dr. Kühler's Conception of Early Dutch Anabaptism Historically Sound? The Historical Discussion of Münster 450 Years Later." Mennonite Quarterly Review 60 (1986): 261-88.
Voolstra, Sjouke. Het Woord is Vlees Geworden: De Melchioritisch-Menniste Incarnatieleer. Kamkpen: Kok, 1982.
Waite, Gary K. "Spiritualizing the Crusade: David Joris in the Context of the Early Reform and Anabaptist Movements in the Netherlands 1524-1543." PhD diss. U., of Waterloo, 1986.
Wray, Frank J. "The Anabaptist Doctrine of the Restitution of the Church." Mennonite Quarterly Review 28 (1954): 186-96.
Zeman, Jarold K. "Anabaptism: A Replay of Medieval Themes or a Prelude to the Modern Age." Mennonite Quarterly Review 50 (1976): 259-71.
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Adapted by permission of Herald Press, Harrisonburg, Virginia, and Waterloo, Ontario, from Mennonite Encyclopedia, Vol. 1, pp. 113-116; vol. 5, pp. 23-26. All rights reserved. For information on ordering the encyclopedia visit the Herald Press website.
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MLA style: Bender, Harold S., Robert Friedmann and Walter Klaassen. "Anabaptism." Global Anabaptist Mennonite Encyclopedia Online. 1990. Web. 19 June 2013. http://www.gameo.org/encyclopedia/contents/A533ME.html.
APA style: Bender, Harold S., Robert Friedmann and Walter Klaassen. (1990). Anabaptism. Global Anabaptist Mennonite Encyclopedia Online. Retrieved 19 June 2013, from http://www.gameo.org/encyclopedia/contents/A533ME.html. |
As the war in Vietnam intensified, the issue of prisoners of war and war crimes gained in significance, both in Vietnam and worldwide. The MACV Staff judge Advocate became deeply involved in this issue, developing legal policy and implementing rules based on the Geneva Conventions, United States policy, and Free World forces interests.
Application of Geneva Conventions to Prisoners of War
The application of the Geneva Conventions of 12 August 1949 to captives in the Vietnam War was complicated by the perplexing legal nature of that conflict. In the classic sense, the conventions presume a declared state of war between two or more sovereign states, each fielding a regular army fighting on a readily identifiable battlefront. Virtually none of these classic conditions existed in the Vietnam conflict.
The United States recognized the sovereignty of South Vietnam, as did some eighty-seven other nations. Indeed, South Vietnam is a member of several special committees of the United Nations, and would have been a member of the United Nations itself had it not been for a Soviet veto in 1957.
The United States has not accorded full diplomatic recognition to North Vietnam, as have some twenty-seven other states. However, the United States has acknowledged North Vietnam's agreement to the Geneva Conventions of 1949, and it has treated North Vietnam as a separate state in the context of Article 12 of the Geneva Prisoner of War Conventions.
Throughout the course of the war the government of North Vietnam was most reluctant to admit to any involvement in South Vietnam, constantly maintaining all of Vietnam to be one country and the Saigon government a puppet regime, beleaguered by indigenous patriots who wished to restore the country to the people. The Republic of Vietnam, while asserting its separation from North Vietnam and its existence as a sovereign state, steadfastly refused to accord to
the Viet Cong any degree of legitimacy, either as a separate political entity or as an agent of Hanoi.
Neither the United States nor North Vietnam issued a declaration of war. South Vietnam declared a state of emergency in 1964 and a state of war in 1965, actions taken primarily to increase its internal powers.
The types of combat forces in the war ranged the full spectrum from the regular divisions of the United States, South Vietnam, and North Vietnam to the Regional Forces and Popular Forces of the government of South Vietnam, the Main Force and Local Force battalions of the Viet Cong, the Civilian Irregular Defense Group of South Vietnam, and the secret self-defense groups of the Viet Cong. The battlefield was nowhere and everywhere, with no identifiable front lines, and no safe rear areas. Fighting occurred over the length and breadth of South Vietnam, on the seas, into Laos and Cambodia, and in the air over North Vietnam. It involved combatants and civilians from a dozen different nations. Politically, militarily, and in terms of international law, the Vietnam conflict posed problems of deep complexity. The inherent difficulty of attempting to apply traditional principles of international law to such a legally confusing conflict is well illustrated by the issue of prisoners of war.
As combat units of the United States became heavily engaged in the war in 1965, the question arose as to the proper disposition for battlefield captives and others detained by U.S. units during military operations. In 1965 the United States determined to win over to the Vietnamese armed forces all individuals captured by U.S. forces. Such an arrangement is permissible under the Geneva Prisoner of War Conventions, which provide for the capturing power to release prisoners to a detaining power as long as both the capturing and the detaining powers fulfill certain obligations concerning the welfare of the prisoners.
While the legal basis for a transfer of prisoners was sound, carrying out the transfer was beset by serious legal and practical difficulties. The Republic of Vietnam regarded the Viet Cong as criminals who violated the security laws of South Vietnam and who consequently were subject to trial for their crimes. As indigenous offenders, the Viet Cong did not technically merit prisoner of war status, although they were entitled to humane treatment under Article 3, Geneva Prisoner of War Conventions. Under Article 12, the United States retained responsibility for treatment of its captives in accordance with the Geneva Conventions even after transfer of the captives to the South Vietnamese. At the same time, the United States was concerned that Americans held captive in North and South Vietnam receive humane treatment and be accorded the full benefits and protection of prisoners of war. In the south, where the government of
South Vietnam had tried and publicly executed some Viet Cong agents, there had been retributory executions of Americans by the Viet Cong. In the north, the Hanoi government stated that it would treat captured American flyers humanely, but it would not accord them prisoner of war status as they were "pirates" engaged in unprovoked attacks on North Vietnam. Hanoi repeatedly threatened to try United States pilots in accordance with Vietnamese laws, but never carried out this threat. U.S. policy was for the United States to do all in its power to alleviate the plight of American prisoners. It was expected that efforts by the United States to ensure humane treatment for Viet Cong and North Vietnamese Army captives would bring reciprocal benefits for American captives.
Early in the war there had been some question in the United States command as to whether the struggle against the Viet Cong constituted an armed international conflict as contemplated in Article 2, Geneva Prisoner of War Conventions, or a conflict not of international nature, to which Article 3 would be applicable. With the infusion of large numbers of United States and North Vietnamese combat units and the coming of the Korean, Australian, Thai, and New Zealand contingents of the Free World Military Assistance Forces, any practical doubts as to the international nature of the conflict were resolved. Although North Vietnam made a strong argument that the conflict in Vietnam was essentially an internal domestic struggle, the official position of the United States, stated as early as 1965, and repeated consistently thereafter, was that the hostilities constituted an armed international conflict, that North Vietnam was a belligerent, that the Viet Cong were agents of the government of North Vietnam, and that the Geneva Conventions applied in full. This view was urged upon the government of South Vietnam, which acceded reluctantly, but subsequently came out in full support of the conventions.
A major practical difficulty in implementing a prisoner of war program was that the Vietnamese government had no facilities suitable for the confinement and care of prisoners of war. In December 1964, the Vietnamese Director of Military justice took the MACV Staff judge Advocate on a tour of courts and confinement facilities throughout South Vietnam. As a result of his observations during that tour the Staff judge Advocate prepared a study pointing out some of the serious problems that existed in handling Viet Cong suspects and prisoners. These problems were quickly becoming joint U.S.-Vietnamese problems, because combat captives and Viet Cong suspects picked up by U.S. forces, Free World Military Assistance Forces, and Vietnamese forces were all delivered to Vietnamese authorities for interrogation, processing, and possible confinement.
During 1965 the number of political prisoners in confinement
rose almost 100 percent, from 9,895 in January to 18,786 in December. These were primarily members of the Viet Cong, but also included some Viet Cong sympathizers, supporters, or collaborators. A total of 24,878 of these political prisoners were confined during the year (compared with 14,029 in 1964), while 15,987 such prisoners were released during the same period. The total rated capacity of all South Vietnam civilian jails and prisons was about 21,400. Few political prisoners, terrorists, or prisoners of war were customarily held in Vietnamese military prisons, which were used primarily as pretrial detention centers for Army deserters and other military offenders. After June 1965 the prison population steadily rose until by early 1966 there was no space for more prisoners in the existing jails and prisons. The practical effect of this was that as new prisoners were confined others had to be discharged. Average time of confinement for all prisoners, including Viet Cong, was about six months. Thus a few months after apprehension a Viet Cong member could be, and usually was, back in operation, while had he been a prisoner of war he would have been restrained "for the duration."
Lack of physical space was but one of many serious problems. An important factor in the operations of the jails and prisons was simply the cost of feeding the prisoners. At an average allowance of 14 piasters (100) per prisoner per day, the monthly cost by early 1966 ran to about ten and a half million piasters a month. Confinement authorities complained of chronic difficulty in feeding the large number of prisoners they-were. required to care for, and the jailer in one province was understandably reluctant to accept and feed prisoners from other provinces. Again, the result was that many prisoners were released after a short time simply because they could not be fed.
Confinement facilities were also handicapped by a severe shortage of qualified administrative and security officers. As an illustration, in 1965 at the Tam Hiep prison there were 30 guards for all shifts and 50 other personnel to control and train almost 1,000 prisoners. In 1966 at Nha Trang Province jail there were 26 guards and 1 instructor for about 440 prisoners. The situation was further aggravated by the frequent loss of guards, jailers, wardens, and instructors who were drafted into the military. This manpower shortage not only thwarted any meaningful classification and rehabilitation program for the prisoners, but also seriously threatened the security of jails and prisons, which were prime Viet Cong targets as long as they held Viet Cong prisoners.
In terms of the war effort, probably the most serious shortcoming of the prisons was the fact that common criminals, Viet Cong suspects, prisoners of war, and even juvenile delinquents were all mixed together. This enabled Viet Cong agents to foment resentment against the government of the Republic of Vietnam and to proselitize their
fellow prisoners; it also increased a Viet Cong's chance for early release as part of the normal inflow-outflow of the prison population. The following figures show the population of Vietnamese jails and prisons that housed political prisoners as of the end of 1965. The figures do not include military prisons, which at that time also held some Viet Cong who were taken in combat and more closely resembled the classical prisoners of war.
|Prisons||Total prisoners||Viet Cong||Prisoners of war|
|Con Son National Prison||3,551||2,934||41|
|Chi Hoa National Prison||4,179||1,513||0|
|Tam Hiep National Prison||946||809||137|
|Thu Duc National Prison||673||296||0|
|42 Province jails||14,035||12,069||886|
Three possible ways of alleviating the overcrowded conditions in the prisons, brought on by the escalation of the war, were suggested: a prisoner of war camp construction program; a broadening of the prisoner of war concept beyond the terms of Article 4 of the Geneva Prisoner of War Conventions so as to include more Viet Cong in the prisoner of war category; and the establishment by the Vietnamese government of re-education centers to separate and rehabilitate suspects who either did not qualify for prisoner of war status or were to be brought before a criminal court as civilian defendants.
In August 1965 the U.S. government and the Vietnamese government notified the International Committee of the Red Cross that their armed forces were abiding by and would continue to abide by the Geneva Conventions. In September a Vietnamese-U.S. joint military committee was appointed to work out details on the application of the Geneva Prisoner of War Conventions in Vietnam. By October the committee had issued three-by-five-inch cards and other training aids for the troops, explaining prisoner of war treatment under the Geneva Conventions. A program of instruction for all U.S. and Vietnamese military units was established to teach the basic rules to be applied in the handling of prisoners. U.S. units and advisers were instructed to identify and record all captives turned over to the Vietnamese, specifying to whom the captives were transferred. Vietnamese military legal advisers were briefed by the MACV Staff judge Advocate on the legal aspects of applying the conventions. The Commander, U.S. Military Assistance Command, Vietnam, established a policy that all suspected Viet Cong captives taken by U.S. forces were to be treated initially as prisoners of war by the capturing unit. Capturing units were responsible for all of the enemy taken prisoner during the course of operations, from the time of their capture to the time the prisoners were released to Vietnamese authorities. Captives were to be interrogated and detained by U.S. forces
only long enough to obtain from them any legitimate tactical intelligence they possessed. Captives were then to be sent to a combined U.S.-Vietnamese Army interrogation center for classification and further processing. Prisoners of war were sent to prisoner of war camps; innocent civilians were released and returned to the place of capture, if possible; civilian defendants were turned over to Vietnamese civil authorities or the province security committee; former Viet Cong seeking amnesty under the Chieu Hoi (Open Arms) program were sent to the Chieu Hoi center. Chieu Hoi was an amnesty program established by the Vietnamese government to encourage Viet Cong to return to government control.
The classification of Viet Cong combatants and Viet Cong suspects posed an interesting legal problem. Because it believed the Viet Cong were traitors and criminals, the Vietnam government was reluctant to accord prisoner of war status to Viet Cong captives. Furthermore it was certainly arguable that many Viet Cong did not meet the criteria of guerrillas entitled to prisoner of war status under Article 4, Geneva Prisoner of War Conventions. However, civil incarceration and criminal trial of the great number of Viet Cong was too much for the civil resources at hand. In addition, Article 22 prohibited the mingling of civil defendants with prisoners of war. By broadly construing Article 4, so as to accord full prisoner of war status to Viet Cong Main Force and Local Force troops, as well as regular North Vietnamese Army troops, any Viet Cong taken in combat would be detained for a prisoner of war camp rather than a civilian jail. The MACV policy was that all combatants captured during military operations were to be accorded prisoner of war status, irrespective of the type of unit to which they belonged. Terrorists, spies, and saboteurs were excluded from consideration as prisoners of war. Suspected Viet Cong captured under circumstances not warranting their treatment as prisoners of war were handled as civilian defendants. MACV policy concerning the classification and treatment of prisoners of war was first codified in MACV Directive 381-11, dated 5 March 1966. (See Appendix D.)
The delegate of the International Committee of the Red Cross speaking in Saigon had the following to say about the MACV policy concerning treatment of Viet Cong as prisoners of war:
The MACV instruction . . . is a brilliant expression of a liberal and realistic attitude.... This text could very well be a most important one in the history of the humanitarian law, for it is the first time . . . that a government goes far beyond the requirements of the Geneva Convention in an official instruction to its armed forces. The dreams of today are the realities of tomorrow, and the day those definitions or similar ones will become embodied in an international treaty . . . will be a great one for man concerned about the protection of men who cannot protect themselves.
... May it then be remembered that this light first shone in the darkness of this tragic war of Vietnam.
Establishing a Prisoner of War Program
On 27 November 1965 the joint Military Committee proposed a workable plan for application of the Geneva Prisoner of War Conventions by the U.S., Vietnamese, and Free World forces. The plan called for the construction of five prisoner of war camps, one in each corps tactical zone and one in the Capital Military Region (Saigon), each having an initial capacity of 1,000 prisoners. Each camp would be staffed by Vietnamese military police, with U.S. military police prisoner of war advisers also assigned to each stockade. The plan was approved in December 1965-a temporary prisoner of war camp was to be established at Bien Hoa in early January 1966, with permanent prisoner of war camps to follow. Prisoner of war camp construction continued to receive priority command attention in 1966. The Bien Hoa camp in III Corps was opened in May, the Pleiku camp in II Corps was completed in August, and the Da Nang camp in I Corps was opened in November. Late in the year work was begun on the Can Tho camp in IV Corps.
The objectives for the prisoner of war program for 1967 were ambitious: to identify and transfer prisoners of war in civilian jails and prisons to Vietnamese Army prisoner of war camps; to establish a program of repatriation of prisoners of war; to promulgate the provisions of the Geneva Conventions of 1949; to establish effective prisoner of war accountability procedures and maintain records for the identification and handling of prisoners of war; to construct additional prisoner of war camps as required; to establish prisoner of war labor and educational programs; to adhere to the Geneva Prisoner of War Conventions as closely as possible with respect to mail, medical attention, Red Cross visits, visiting privileges, and health and welfare.
By the end of 1967 the prisoner of war camp capacity had increased from 3,000 to 13,000. In March 1968 a camp for female prisoners of war was established at Qui Nhon, and in April steps were taken to concentrate all Viet Cong prisoners of war under age eighteen at the Bien Hoa camp, where they received special rehabilitation, education, and vocational training. A central prisoner of war camp was constructed at Phu Quoc Island, off the coast of Cambodia, and by the end of 1968 the prisoner of war camps could house a population of 21,000 normally, and a total of 32,000 on an emergency or short-range basis. All gradually expanded until by 11 December 1971 the Vietnamese government held 35,665 prisoners of war in six camps. Of these, 13,365 had been captured by U.S. forces.
Inspections by the International Committee of the Red Cross
Construction of the prisoner camps was a major feat in itself, but the U.S. and Vietnamese governments worked hard in many other areas to fulfill their responsibilities under the Geneva Conventions. At first the Vietnam government was reluctant to co-operate with the International Committee of the Red Cross with respect to inspections and furnishing lists of prisoners. Furthermore, the Democratic Republic of North Vietnam refused to allow the Red Cross any access to their prisoners and the South Vietnamese felt there should be reciprocity. In South Vietnam confinement facilities were the responsibility of the Minister of the Interior, and at the urging of the United States he agreed to allow more visits by Red Cross representatives to Vietnamese civil prisons and re-education centers where prisoners of war were being detained until the completion of the camps. In early 1966, as a result of U.S. efforts, representatives of the International Committee of the Red Cross visited prisons at Tam Hiep, Con Son, and Da Nang, and the prisoner of war camp under construction at Bien Hoa. After it was completed, the Bien Hoa camp was again visited by Red Cross representatives in August 1966. The representatives were favorably impressed with the camp and agreed to provide health and welfare items on their next visit. In October 1966 the committee's representatives visited detention facilities in Da Nang and Pleiku, and again they were favorably impressed. Another Red Cross representative, accompanied by a Saigon delegation, made an extensive tour from 29 November until 8 December 1966 of South Vietnamese, Free World, and U.S. prisoner of war facilities throughout the Republic of Vietnam. The representative visited Vietnamese prisoner of war camps in I, II, and III Corps Tactical Zones; two Vietnamese Army hospitals; the Australian prisoner of war collecting point and field hospital; the Republic of Korea Capitol Division collecting point and hospital; the four U.S. collecting points and six U.S. hospitals; the III Marine Amphibious Force special detention facility at Da Nang; and the Philippine hospital at Tay Ninh. In 1967 members of the international press accompanied the representatives on their visits to two camps.
Despite the many problems they encountered, the record is clear that the United States and Vietnam made a vigorous effort to adhere to the exacting standards of the Geneva Prisoner of War Conventions. Within the Military Assistance Command, Vietnam, the provost marshal was responsible for advising the Vietnamese prisoner of war camps, ensuring that they were operated in conformity with Geneva requirements, and acting as the point of contact for representatives of the International Committee of the Red Cross. MACV policy and procedures pertaining to Red Cross inspections of prisoner of war
INTERNATIONAL COMMITTEE OF THE RED CROSS VISITS VIETNAMESE ARMY PRISONER OF WAR CAMP AT LONG BINH
camps were set forth in MACV Directive 190-6. (See Appendix E.)
In August 1966 a committee began to screen prisoners who had been placed in civilian jails before the construction of the prisoner of war camps. The committee, which was composed of representatives from MACV directorates of personnel, intelligence, and logistics and their Vietnamese counterparts, screened all four national prisons and 37 provincial jails. The committee identified 1,202 men as prisoners of war; all but 27 of these had been transferred to prisoner of war camps by the end of 1967, as the screening process continued.
Several programs to improve the living standards of prisoners of war were approved in 1967, among them a pilot educational program to teach reading and writing. Mailing privileges were granted the prisoners and visitation rights given the families of Viet Cong prisoners of war. Health and comfort items were issued free, and a dispensary with a medical doctor and staff was established at each camp. In August 1967 the MACV directorate of personnel established a Prisoner of. War Work Advisory Detachment to encourage and assist work programs at the prisoner of war camps. In December 1967, the Vietnamese government approved a labor program whereby prisoners of war would work and be paid 8 piasters per day. This program, which was in accordance with recommendations of the Geneva Prisoner of War Conventions, became effective 1 April 1968.
Throughout the course of the war the allies maintained a major and sustained effort to promote a reciprocal program of prisoner repatriation. On 3 February 1967, twenty-eight North Vietnamese prisoners of war were released at the Ben Hai River to return to North Vietnam through the demilitarized zone. On 11 March 1967, two Viet Cong prisoners of war captured by U.S. forces were released in response to the release of two U.S. prisoners of war. On 20 March two North Vietnamese PT boat crewmen were repatriated through Cambodia; on 22 March the South Vietnamese released twenty-two Viet Cong prisoners of war and on 23 June 1967 three Viet Cong captured by U.S. forces were released at a jungle rendezvous in exchange for the release of two U.S. prisoners and one Filipino captured by the Viet Cong. In April 1967 a screening program was begun to identify prisoners of war who, because of illness, were qualified for release under Articles 109 and 110 of the Geneva Conventions. The screening team included two Swiss physicians under contract to the International Committee of the Red Cross (ICRC). Of the 286 prisoners screened, 135 qualified medically for repatriation. Of those qualified for repatriation, only 39 wished to return to North Vietnam. To this group was added a female prisoner of war who had given birth in a South Vietnamese hospital. The 40 prisoners and the infant were repatriated to North Vietnam through the demilitarized zone on 12 June 1967; on the same day four Viet Cong-U.S. prisoners were released in South Vietnam. During 1967 a total of 139 prisoners of war were released in South Vietnam or repatriated to North Vietnam.
In 1968 the government of South Vietnam, with U.S. support, sought to repatriate 40 sick and wounded prisoners of war to North Vietnam under Articles 109 and 110 of the Geneva Conventions. The prisoners were examined by a Red Cross physician and had expressed a desire to return to the north. The Vietnamese made the repatriation offer through the Red Cross, which sent a telegram to the North Vietnamese Minister of Foreign Affairs, proposing repatriation in late January or early February. When no reply was received, the South Vietnamese asked the International Committee of the Red Cross to renew its efforts in May, but there was still no reply by the end of the year. Efforts by the South Vietnamese and the Red Cross to repatriate these 40 prisoners of war, and 24 civilians as well, continued through 1969, but to no avail. At the Paris Peace Talks of 13 November 1969 the South Vietnamese proposed returning 62 sick and wounded prisoners of war to North Vietnam. The offer was declined. It was the position of the South Vietnamese that if Geneva Prisoner of War Articles 109 and 110 required the capturing state to
repatriate the sick and wounded, these articles also required the home state of the prisoners of war to accept those prisoners who wished repatriation. The South Vietnamese even proposed transporting sick and wounded prisoners of war by sea to any port or point on the coast of North Vietnam, but received no response to this offer.
During 1969 the Republic of Vietnam did release 191 Viet Cong prisoners of war in South Vietnam for reasons of youth, age, or pregnancy under Article 21 of the Geneva Prisoner of War Conventions. Also, some prisoners of war who were not deemed to be "hard-core" Viet Cong were transferred to the favored Hoi Chanh status, under the Chieu Hoi amnesty program. These transfers particularly applied to youths seventeen and under, Efforts to repatriate prisoners of war to North Vietnam and to secure the release of U.S. and South Vietnamese prisoners of war continued through 1970, but met with little success. In January 1971 the Republic of Vietnam offered to repatriate all sick and wounded prisoners to North Vietnam. This offer was ignored. On 29 April 1971 the Vietnam government requested North Vietnam to conclude a bilateral agreement for the repatriation or internment in a neutral country of those prisoners of war who had been held captive for a long period of time. This offer was ignored, but in May 1971 North Vietnam finally agreed to accept 570 sick and wounded prisoners. The International Committee of the Red Cross interviewed 660 sick and wounded prisoners, of whom only 13 wished to be repatriated. As arranged, these thirteen were taken by sea to a point off the coast of North Vietnam; but before they were released North Vietnam canceled the agreement to accept the prisoners and they were returned to Da Nang. This ended repatriations for the remainder of 1971.
Throughout the war the United States had urged the South Vietnamese to release qualified prisoners of war, seeking possible reciprocal action by North Vietnam. The South Vietnamese had been understandably reluctant to release unilaterally large numbers of able-bodied prisoners of war, but after the national election of October 1970, the Republic of Vietnam transferred 2,300 Viet Cong prisoners who pledged loyalty to the government to Chieu Hoi centers and released 623 outright. Through I March 1972, South Vietnam released a total of 5,960 prisoners of war. Of this total 188 were repatriated to North Vietnam, 900 were released in South Vietnam, 1,784 were reclassified, and 3,084 were accepted into the Chieu Hoi program. In contrast, by the end of 1971 the Communists had released 53 American prisoners of war; they had allowed no visits by the International Committee of the Red Cross to Communist prisoner of war camps in North Vietnam or South Vietnam; and they had made no effort to repatriate sick and wounded prisoners. Mailing privileges for U.S. prisoners held in North Vietnam had been al-
lowed sporadically and in an arbitrary manner. Finally, the Communists had refused to furnish a comprehensive and accurate list of the prisoners they held. Although both the Viet Cong and North Vietnamese consistently maintained that their prisoners received humane treatment, their efforts to comply with the provisions of the Geneva Prisoner of War Conventions were negligible.
MACV Policy on Prevention and Investigation of War Crimes
The first MACV directive dealing specifically with war crimes was Directive 20-4, dated 20 April 1965. Its purpose was to designate the agencies responsible for conducting investigations of alleged or apparent violations of the Geneva Conventions inflicted by hostile forces upon U.S. military or civilian personnel assigned to Vietnam. The directive defined war crimes as violations of the law of war, and stated that "grave breaches" of the Geneva Conventions, such as willful killing, torture, or inhuman treatment of persons protected by the conventions constituted war crimes. The directive also addressed itself to "prohibited acts" under common Article 3 of the Geneva Conventions, defining them as those acts which "would be war crimes but for the fact the existing conflict is not yet deemed to be international in character." In retrospect, this early directive is interesting in two respects. The directive concerned itself only with those violations inflicted by hostile forces upon U.S. citizens, and it stated that the fighting in Vietnam was not yet an armed international conflict in the legal sense.
The reasons for drafting the directive in this manner were immediate and practical. Prior to the introduction of ground combat units in Vietnam in March 1965, U.S. troops in Vietnam served in an advisory capacity; U.S. units had not planned or executed combat operations or taken prisoners, and there were no indications that U.S. advisers were violating the Geneva Conventions. To the contrary, the only atrocities known to the U.S. command at the time were those committed against U.S. advisers by the Communists. MACV Directive 204 was promulgated to ensure appropriate investigation of such atrocities. (See Appendix F.) As to the international nature of the conflict, the principal enemy of the Republic of Vietnam at the time was the Viet Cong, whose members the republic regarded as domestic criminals. The government of North Vietnam did not admit that its troops were in the south, or that it was in any way sponsoring the Viet Cong, whom it considered local patriots struggling against dictatorial regime.
MACV Directive 20-4 was updated on 25 March 1966 with several significant changes. The scope of the directive was broadened to include war crimes committed by U.S. personnel as well as those against U.S. personnel. In consonance with the official U.S. position
that the struggle in Vietnam by then constituted an armed international conflict, no mention was made of a lesser category of "prohibited acts," as defined in the earlier directive, and eighteen examples of acts which constituted war crimes were given.
The directive clearly stated that the willful killing, torture, or inhuman treatment of, or willfully causing great suffering or serious injury to the body or health of persons taking no active part in the hostilities, including members of the armed forces who had laid down their arms or who were not combatants because of sickness, wounds, or any other cause, was a war crime. Other acts specified as war crimes were maltreatment of dead bodies, firing on localities which were undefended and without military significance, pillage or purposeless destruction, killing without trial of spies or other persons who committed hostile acts, and compelling prisoners of war or civilians to perform labor prohibited by the Geneva Conventions.
The directive went on to fix responsibility on every member of the U.S. armed forces for reporting incidents which could constitute war crimes. It stated that "It is the responsibility of all military personnel having knowledge or receiving a report of an incident or of an act thought to be a war crime to make such incident known to his commanding officer as soon as practicable. Personnel performing investigative, intelligence, police, photographic, grave registration or medical functions, as well as those in contact with the enemy will, in the normal course of their duties, make every effort to detect the commission of war crimes and will report the essential facts to their commanding officer."
As was the-practice for all MACV directives, MACV Directive 20-4 was updated periodically. It was also supplemented by a number of other directives pertaining to the Geneva Conventions, war crimes, and prisoners of war. MACV Directive 27-5, Legal Services: War Crimes, and Other Prohibited Acts, dated 2 November 1967, listed acts which constituted war crimes and stated, "Commission of any act, enumerated in paragraph 4, above, or constituting a war crime is prohibited. Violation of this directive will be punishable in accordance with the provisions of the Uniform Code of Military justice." Other regulations pertinent to war crimes and prisoners of war were MACV Directive 190-3, Military Police: Enemy Prisoners of War, 6 April 1967; MACV Directive 20-5, Inspections and Investigations: Prisoners of War-Determination of Eligibility, 17 May 1966; MACV Directive 381-46, Military Intelligence: Combined Screening of Detainees, 27 December 1967; MACV Directive 335-1, Reports of Serious Crimes or Other Incidents, 5 January 1966; and Geneva Conventions Checksheet sent to all judge advocates in Vietnam. (See Appendix G.)
Throughout 1965, 1966, and 1967 the most grievous breaches of
the Geneva Conventions continued to be those committed by the Communists, and there were several cases where U.S. troops were murdered and their bodies mutilated by the Viet Cong or North Vietnamese. The Viet Cong policy of kidnapping civilians, assassinating public officials, and terrorizing entire populations continued. Communist tactics against the Montagnards, indigenous mountain tribes, were particularly vicious.
On the American side, the sudden massive U.S. troop buildup in Vietnam that began in 1965 created many problems for the U.S. command, and incidents of war crimes by U.S. troops began to be reported. For example, during the period between 1 January 1965 and 31 August 1973, there were 241 cases (excluding My Lai) which involved allegations of war crimes against United States Army troops. One hundred and sixty of these cases, upon investigation, were determined to be unsubstantiated. Substantiated allegations of war crimes Violations committed in Vietnam by personnel subject to the Uniform Code of Military justice were prosecuted under the provisions of the code. From January 1965 through August 1973, 36 cases involving war crimes allegations against Army personnel were tried by court-martial. Sixteen cases involving thirty men resulted in acquittal or dismissal after arraignment. Twenty cases involving thirty-one Army servicemen resulted in conviction. By the time the U.S. troop buildup was in full swing, various MACV directives contained a sufficient body of law clearly to define, prohibit, and provide for the investigation of war crimes. The constant rotation of troops created a continual need to get the information to the troops.
Long before U.S. troop units were engaged in combat in Vietnam, the Army had included in its training programs material designed to inculcate in the troops a knowledge of their rights and obligations under the Geneva Conventions of 1949. Army Regulation 350-216, dated 19 December 1965, placed upon the Commanding General, Continental Army Command, the responsibility of incorporating within appropriate training programs periods of instruction designed to insure that all members of the Army were familiar with the 1949 Geneva Conventions. The soldier's first introduction to the Geneva Conventions was during basic training, where he received two hours of formal instruction, followed by a test, the results of which were noted on his record. During advanced individual training, instruction in the Geneva Conventions was integrated with other subjects and principles were applied during field exercises. Annex B of Continental Army Command Regulation 350-11 dated 15 June 1965, required commanders to take continuing action to incorporate instruction on
the Geneva Conventions in their training programs. Training programs pertaining to the Geneva Conventions were the subject of a comprehensive review, and several steps were taken to increase and improve instruction by the use of training films and combined judge advocate-combat arms officer training teams.
In Vietnam, Geneva Conventions training was intensified and became more formalized as troop strength increased. As early as August 1965, the Commander, U.S. Military Assistance Command, Vietnam, directed that the educational program for all U.S. military personnel in South Vietnam include the issuance of a three-by-five-inch card containing the basic requirements of the Geneva Conventions pertaining to the treatment of prisoners of war. By October 1965 cards had been prepared in English for U.S. personnel and in Vietnamese for Vietnamese armed forces personnel. U.S. units were instructed in the basic rules of handling prisoners and MACV judge advocate personnel briefed Vietnamese military legal personnel on the application of the conventions. Troops arriving in Vietnam received Geneva Conventions orientation during their initial in-processing period, where they also received a copy of the card, "The Enemy in Your Hands." (See Appendix H.)
The Commander, U.S. Military Assistance Command, Vietnam, continuously and emphatically stressed the importance of all troops acting in accordance with the laws of war. In November 1965, he discussed with the chief of staff of the Vietnamese joint General Staff the importance of adhering to the Geneva Conventions pertaining to the treatment of prisoners of war. The commander conferred periodically with Vietnamese officials on this subject and on the importance of proper deportment by all troops in general. The results of this concern became noticeable in several areas, one noteworthy example being the promulgation in March 1967 by the Vietnamese government of a National Decree stating the provisions of the four Geneva Conventions of 12 August 1949. In August 1966 the Commander, U.S. Military Assistance Command, Vietnam, personally wrote separate letters to all of his major commanders on this same theme, stating in part, "Active command interest in this program, in coordination with Republic of Vietnam Armed Forces authorities, which assures that prisoners of war and combat captives are properly processed and handled in accordance with International Law is vital."
This theme was repeated over and over again. A MACV command information bulletin, titled Application of the Geneva Prisoner of War Conventions in Vietnam, dated October 1966, instructed the U.S. troops that the Geneva Prisoner of War Conventions applied to Vietnam even though there was no formal declaration of war by the United States. Moreover, the United States was applying not
only the letter of the law, but also the spirit of the Geneva Conventions, which were designed to protect the individual who could no longer protect himself. Prisoner of war treatment was to be extended to all Viet Cong and to all members of regular North Vietnamese units, whether captured in combat or not, as long as they were not criminals, spies, saboteurs, or terrorists. Criminals, spies, saboteurs, and terrorists were to be given humane treatment and turned over to the Vietnamese government for trial.
The bulletin explained the steps to be taken immediately upon capture of enemy personnel, and stressed that prisoners must be protected from torture, humiliation, degrading treatment, reprisals, or any act of violence. The categories of detained persons (innocent civilians, prisoners of war, returnees, and civil defendants) were listed, with a reference to the MACV directive which outlined processing procedures for each type of detainee. The bulletin went on to explain the importance of observing humanitarian principles in waging war, giving specific reasons why it was in the interest of the U.S. for American troops to treat prisoners humanely. In conclusion, the bulletin urged the troops to follow the rules on the card "The Enemy in Your Hands." In addition, through U.S. judge advocate resources in Vietnam the MACV Staff judge Advocate's office monitored the troop education program and disciplinary aspects of Geneva Conventions violations.
These efforts on the part of the U.S. command were commended by the Red Cross in a letter of 5 January 1968 to W. Averell Harriman, U.S. Ambassador at Large. Samuel A. Gonard, President of the International Committee of the Red Cross, wrote: "We are convinced that in the context of the war in Vietnam the U.S. Forces are devoting a major effort to the spread of knowledge of the Geneva Conventions."
War Crimes Investigation
For the most part, war crimes committed by U.S. forces in Vietnam fell into two principal categories: willful murder or assault of noncombatants; and mutilation and maltreatment of dead bodies. Serious incidents involving assault, rape, and murder that were not directly connected with military operations in the field were not characterized as war crimes but were reported through military police channels as violations of the Uniform Code of Military justice.
Acts constituting war crimes were also offenses against the Uniform Code of Military justice, and as such were investigated by agents of the Criminal Investigation Division. Pertinent MACV directives required a concurrent investigation of war crimes by an investigating officer who was concerned not only with the details
of the crime, such as the persons involved and where, when, and what occurred, but also with the broader question of how and why the incident took place. The scope of this investigation included an examination of the established rules of engagement and command and control procedures that were in effect at the time, and how these procedures were implemented. The question to be determined was whether there was any failure of command responsibility.
When an investigation was completed, the report was delivered to the general court-martial convening authority, who had appointed the investigating officer. The appointing authority reviewed the report and approved or disapproved it. If approved, the report of the investigation with the appointing authority's indorsement was forwarded through channels to the Commander, U.S. Military Assistance Command, Vietnam. At MACV headquarters it was circulated to appropriate staff offices, including the Staff judge Advocate, for review. The report could be returned for further action or approved by the MACV commander or chief of staff. After final review, a war crimes investigation report concerning any person was forwarded to The judge Advocate General, Department of the Army.
The Commander, U.S. Military Assistance Command, Vietnam, had considered establishing special war crimes teams and having the Army maintain centralized files on war crimes for all services, but this was not done because the laws prohibiting war crimes and the administrative and judicial machinery for investigating and punishing such offenses were judged adequate. Murder, rape, assault, arson, pillage, and larceny were all punishable as offenses against various sections of the Uniform Code of Military justice, and there were many directives from Military Assistance Command, U.S. Army, Vietnam, and units specifying and prohibiting various acts in the war crimes category. Representatives of the military police, Criminal Investigation Division, Inspector General, and judge Advocate had experience in conducting investigations; they, as well as the commanders, and, indeed, all military personnel, had the responsibility for reporting possible violations of the laws of war so that an appropriate investigation could be conducted as specified by regulation.
Despite laws and preventive education, war crimes were committed. Most were isolated incidents, offenses committed by individual U.S. soldiers or small groups. Investigations were conducted, and the records of courts-martial proceedings contain the cases of individuals who were tried and punished. My Lai, the most notorious offense committed by U.S. troops in combat in Vietnam, was not the result of inadequate laws or lack of command emphasis on those laws; it was the failure of unit leaders to enforce the clear
and well-known procedures set forth in applicable regulations. It is tragically true that troops on both sides committed atrocities; but had it not been for the genuine concern of commanders at the highest levels that U.S., Vietnamese, and allied forces conduct themselves humanely and in accordance with the laws of war, the Vietnam War probably would have been far more brutal.
It was evident that international law was inadequate to protect victims in wars of insurgency and counterinsurgency, civil war, and undeclared war. The efforts of the international community to codify the humanitarian law of war in 1949 drew upon examples from World War II which simply did not fit in Vietnam. The law left much room for expediency, political manipulation, and propaganda. The hazy line between civilian and combatant became even vaguer in Vietnam.
There was an absence of effective power to insure compliance by both sides with the Geneva Conventions or to give reassurance of at least minimum protection for victims of the armed conflict. The inability of the International Committee of the Red Cross to function effectively on behalf of U.S., South Vietnamese, or other Free World forces was particularly tragic. The law of war was not completely ineffective; certainly many combatants and noncombatants survived the war because of the application of the law of war. The U.S. military lawyers' role in applying the known and developed humanitarian rules for armed conflict brought credit to the legal profession and to the U.S. Armed Forces.
page updated 30 May 2001
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Social Network Analysis diagram
|Theory · History|
Graph · Complex network · Contagion
|Types of Networks|
|Metrics and Algorithms|
A social network is a social structure made up of a set of social actors (such as individuals or organizations) and a complex set of the dyadic ties between these actors. The social network perspective provides a clear way of analyzing the structure of whole social entities. The study of these structures uses social network analysis to identify local and global patterns, locate influential entities, and examine network dynamics.
Social networks and the analysis of them is an inherently interdisciplinary academic field which emerged from social psychology, sociology, statistics, and graph theory. Georg Simmel authored early structural theories in sociology emphasizing the dynamics of triads and "web of group affiliations." Jacob Moreno is credited with developing the first sociograms in the 1930s to study interpersonal relationships. These approaches were mathematically formalized in the 1950s and theories and methods of social networks became pervasive in the social and behavioral sciences by the 1980s. Social network analysis is now one of the major paradigms in contemporary sociology, and is also employed in a number of other social and formal sciences. Together with other complex networks, it forms part of the nascent field of network science.
A social network is a theoretical construct useful in the social sciences to study relationships between individuals, groups, organizations, or even entire societies (social units, see differentiation). The term is used to describe a social structure determined by such interactions. The ties through which any given social unit connects represent the convergence of the various social contacts of that unit. This theoretical approach is, necessarily, relational. An axiom of the social network approach to understanding social interaction is that social phenomena should be primarily conceived and investigated through the properties of relations between and within units, instead of the properties of these units themselves. Thus, one common criticism of social network theory is that individual agency is often ignored although this may not be the case in practice (see agent-based modeling). Precisely because many different types of relations, singular or in combination, form these network configurations, network analytics are useful to a broad range of research enterprises. In social science, these fields of study include, but are not limited to anthropology, biology, communication studies, economics, geography, information science, organizational studies, social psychology, sociology, and sociolinguistics.
In the late 1800s, both Émile Durkheim and Ferdinand Tönnies foreshadowed the idea of social networks in their theories and research of social groups. Tönnies argued that social groups can exist as personal and direct social ties that either link individuals who share values and belief (Gemeinschaft, German, commonly translated as "community") or impersonal, formal, and instrumental social links (Gesellschaft, German, commonly translated as "society"). Durkheim gave a non-individualistic explanation of social facts, arguing that social phenomena arise when interacting individuals constitute a reality that can no longer be accounted for in terms of the properties of individual actors. Georg Simmel, writing at the turn of the twentieth century, pointed to the nature of networks and the effect of network size on interaction and examined the likelihood of interaction in loosely-knit networks rather than groups.
Major developments in the field can be seen in the 1930s by several groups in psychology, anthropology, and mathematics working independently. In psychology, in the 1930s, Jacob L. Moreno began systematic recording and analysis of social interaction in small groups, especially classrooms and work groups (see sociometry). In anthropology, the foundation for social network theory is the theoretical and ethnographic work of Bronislaw Malinowski, Alfred Radcliffe-Brown, and Claude Lévi-Strauss. A group of social anthropologists associated with Max Gluckman and the Manchester School, including John A. Barnes, J. Clyde Mitchell and Elizabeth Bott Spillius, often are credited with performing some of the first fieldwork from which network analyses were performed, investigating community networks in southern Africa, India and the United Kingdom. Concomitantly, British anthropologist S.F. Nadel codified a theory of social structure that was influential in later network analysis. In sociology, the early (1930s) work of Talcott Parsons set the stage for taking a relational approach to understanding social structure. Later, drawing upon Parsons' theory, the work of sociologist Peter Blau provides a strong impetus for analyzing the relational ties of social units with his work on social exchange theory. By the 1970s, a growing number of scholars worked to combine the different tracks and traditions. One group consisted of sociologist Harrison White and his students at the Harvard University Department of Social Relations. Also independently active in the Harvard Social Relations department at the time were Charles Tilly, who focused on networks in political and community sociology and social movements, and Stanley Milgram, who developed the "six degrees of separation" thesis.Mark Granovetter and Barry Wellman are among the former students of White who elaborated and championed the analysis of social networks.
Levels of analysis
In general, social networks are self-organizing, emergent, and complex, such that a globally coherent pattern appears from the local interaction of the elements that make up the system. These patterns become more apparent as network size increases. However, a global network analysis of, for example, all interpersonal relationships in the world is not feasible and is likely to contain so much information as to be uninformative. Practical limitations of computing power, ethics and participant recruitment and payment also limit the scope of a social network analysis. The nuances of a local system may be lost in a large network analysis, hence the quality of information may be more important than its scale for understanding network properties. Thus, social networks are analyzed at the scale relevant to the researcher's theoretical question. Although levels of analysis are not necessarily mutually exclusive, there are three general levels into which networks may fall: micro-level, meso-level, and macro-level.
|This section requires expansion with: additional examples and references for each sub-level. (January 2012)|
Micro level
At the micro-level, social network research typically begins with an individual, snowballing as social relationships are traced, or may begin with a small group of individuals in a particular social context.
Dyadic level: A dyad is a social relationship between two individuals. Network research on dyads may concentrate on structure of the relationship (e.g. multiplexity, strength), social equality, and tendencies toward reciprocity/mutuality.
Triadic level: Add one individual to a dyad, and you have a triad. Research at this level may concentrate on factors such as balance and transitivity, as well as social equality and tendencies toward reciprocity/mutuality.
Actor level: The smallest unit of analysis in a social network is an individual in their social setting, i.e., an "actor" or "ego". Egonetwork analysis focuses on network characteristics such as size, relationship strength, density, centrality, prestige and roles such as isolates, liaisons, and bridges. Such analyses, are most commonly used in the fields of psychology or social psychology, ethnographic kinship analysis or other genealogical studies of relationships between individuals.
Subset level: Subset levels of network research problems begin at the micro-level, but may cross over into the meso-level of analysis. Subset level research may focus on distance and reachability, cliques, cohesive subgroups, or other group actions or behavior.
Meso level
In general, meso-level theories begin with a population size that falls between the micro- and macro-levels. However, meso-level may also refer to analyses that are specifically designed to reveal connections between micro- and macro-levels. Meso-level networks are low density and may exhibit causal processes distinct from interpersonal micro-level networks.
Organizations: Formal organizations are social groups that distribute tasks for a collective goal. Network research on organizations may focus on either intra-organizational or inter-organizational ties in terms of formal or informal relationships. Intra-organizational networks themselves often contain multiple levels of analysis, especially in larger organizations with multiple branches, franchises or semi-autonomous departments. In these cases, research is often conducted at a workgroup level and organization level, focusing on the interplay between the two structures.
Randomly-distributed networks: Exponential random graph models of social networks became state-of-the-art methods of social network analysis in the 1980s. This framework has the capacity to represent social-structural effects commonly observed in many human social networks, including general degree-based structural effects commonly observed in many human social networks as well as reciprocity and transitivity, and at the node-level, homophily and attribute-based activity and popularity effects, as derived from explicit hypotheses about dependencies among network ties. Parameters are given in terms of the prevalence of small subgraph configurations in the network and can be interpreted as describing the combinations of local social processes from which a given network emerges. These probability models for networks on a given set of actors allow generalization beyond the restrictive dyadic independence assumption of micro-networks, allowing models to be built from theoretical structural foundations of social behavior.
Scale-free networks: A scale-free network is a network whose degree distribution follows a power law, at least asymptotically. In network theory a scale-free ideal network is a random network with a degree distribution that unravels the size distribution of social groups. Specific characteristics of scale-free networks vary with the theories and analytical tools used to create them, however, in general, scale-free networks have some common characteristics. One notable characteristic in a scale-free network is the relative commonness of vertices with a degree that greatly exceeds the average. The highest-degree nodes are often called "hubs", and may serve specific purposes in their networks, although this depends greatly on the social context. Another general characteristic of scale-free networks is the clustering coefficient distribution, which decreases as the node degree increases. This distribution also follows a power law. The Barabási model of network evolution shown above is an example of a scale-free network.
Macro level
Large-scale networks: Large-scale network is a term somewhat synonymous with "macro-level" as used, primarily, in social and behavioral sciences, in economics. Originally, the term was used extensively in the computer sciences (see large-scale network mapping).
Complex networks: Most larger social networks display features of social complexity, which involves substantial non-trivial features of network topology, with patterns of complex connections between elements that are neither purely regular nor purely random (see, complexity science, dynamical system and chaos theory), as do biological, and technological networks. Such complex network features include a heavy tail in the degree distribution, a high clustering coefficient, assortativity or disassortativity among vertices, community structure, and hierarchical structure. In the case of agency-directed networks these features also include reciprocity, triad significance profile (TSP, see network motif), and other features. In contrast, many of the mathematical models of networks that have been studied in the past, such as lattices and random graphs, do not show these features.
Imported theories
Various theoretical frameworks have been imported for the use of social network analysis. The most prominent of these are Graph theory, Balance theory, Social comparison theory, and more recently, the Social identity approach.
Indigenous theories
Few complete theories have been produced from social network analysis. Two that have are Structural Role Theory and Heterophily Theory.
The basis of Heterophily Theory was the finding in one study that more numerous weak ties can be important in seeking information and innovation, as cliques have a tendency to have more homogeneous opinions as well as share many common traits. This homophilic tendency was the reason for the members of the cliques to be attracted together in the first place. However, being similar, each member of the clique would also know more or less what the other members knew. To find new information or insights, members of the clique will have to look beyond the clique to its other friends and acquaintances. This is what Granovetter called "the strength of weak ties."
Structural holes
In the context of networks, social capital exists where people have an advantage because of their location in a network. Contacts in a network provide information, opportunities and perspectives that can be beneficial to the central player in the network. Most social structures tend to be characterized by dense clusters of strong connections. Information within these clusters tends to be rather homogeneous and redundant. Non-redundant information is most often obtained through contacts in different clusters. When two separate clusters possess non-redundant information, there is said to be a structural hole between them. Thus, a network that bridges structural holes will provide network benefits that are in some degree additive, rather than overlapping. An ideal network structure has a vine and cluster structure, providing access to many different clusters and structural holes.
Information benefits
Networks rich in structural holes are a form of social capital in that they offer information benefits. The main player in a network that bridges structural holes is able to access information from diverse sources and clusters. This is beneficial to an individual’s career because he is more likely to hear of job openings and opportunities if his network spans a wide range of contacts in different industries/sectors. This concept is similar to Mark Granovetter’s theory of weak ties, which rests on the basis that having a broad range of contacts is most effective for job attainment.
Social capital mobility benefits
In many organizations, members tend to focus their activities inside their own groups, which stifles creativity and restricts opportunities. A player whose network bridges structural holes has an advantage in detecting and developing rewarding opportunities. Such a player can mobilize social capital by acting as a “broker” of information between two clusters that otherwise would not have been in contact, thus providing access to new ideas, opinions and opportunities. British philosopher and political economist John Stuart Mill, writes, “it is hardly possible to overrate the value...of placing human beings in contact with persons dissimilar to themselves…Such communication [is] one of the primary sources of progress.” Thus, a player with a network rich in structural holes can add value to an organization through new ideas and opportunities. This in turn, helps an individual’s career development and advancement.
A social capital broker also reaps control benefits of being the facilitator of information flow between contacts. In the case of consulting firm Eden McCallum, the founders were able to advance their careers by bridging their connections with former big 3 consulting firm consultants and mid-size industry firms. By bridging structural holes and mobilizing social capital, players can advance their careers by executing new opportunities between contacts.
There has been research that both substantiates and refutes the benefits of information brokerage. A study of high tech Chinese firms by Zhixing Xiao found that the control benefits of structural holes are “dissonant to the dominant firm-wide spirit of cooperation and the information benefits cannot materialize due to the communal sharing values” of such organizations. However, this study only analyzed Chinese firms, which tend to have strong communal sharing values. Information and control benefits of structural holes are still valuable in firms that are not quite as inclusive and cooperative on the firm-wide level. In 2004, Ronald Burt studied 673 managers who ran the supply chain for one of America’s largest electronics companies. He found that managers who often discussed issues with other groups were better paid, received more positive job evaluations and were more likely to be promoted. Thus, bridging structural holes can be beneficial to an organization, and in turn, to an individual’s career.
Research clusters
|This section requires expansion with: additional theoretical perspectives and additional examples and references for existing areas of theory. (January 2012)|
Communication Studies are often considered a part of both the social sciences and the humanities, drawing heavily on fields such as sociology, psychology, anthropology, information science, biology, political science, and economics as well as rhetoric, literary studies, and semiotics. Many communications concepts describe the transfer of information from one source to another, and can thus be conceived of in terms of a network.
In J.A. Barnes' day, a "community" referred to a specific geographic location and studies of community ties had to do with who talked, associated, traded, and attended church with whom. Today, however, there are extended "online" communities developed through telecommunications devices and social network services. Such devices and services require extensive and ongoing maintenance and analysis, often using network science methods. Community development studies, today, also make extensive use of such methods.
Complex networks
Criminal networks
In criminology and urban sociology, much attention has been paid to the social networks among criminal actors. For example, Andrew Papachristos has studied gang murders as a series of exchanges between gangs. Murders can be seen to diffuse outwards from a single source, because weaker gangs cannot afford to kill members of stronger gangs in retaliation, but must commit other violent acts to maintain their reputation for strength.
Diffusion of innovations
Diffusion of ideas and innovations studies focus on the spread and use of ideas from one actor to another or one culture and another. This line of research seeks to explain why some become "early adopters" of ideas and innovations, and links social network structure with facilitating or impeding the spread of an innovation.
In demography, the study of social networks has led to new sampling methods for estimating and reaching populations that are hard to enumerate (for example, homeless people or intravenous drug users.) For example, respondent driven sampling is a network-based sampling technique that relies on respondents to a survey recommending further respondents.
Economic sociology
The field of sociology focuses almost entirely on networks of outcomes of social interactions. More narrowly, economic sociology considers behavioral interactions of individuals and groups through social capital and social "markets". Sociologists, such as Mark Granovetter, have developed core principles about the interactions of social structure, information, ability to punish or reward, and trust that frequently recur in their analyses of political, economic and other institutions. Granovetter examines how social structures and social networks can affect economic outcomes like hiring, price, productivity and innovation and describes sociologists’ contributions to analyzing the impact of social structure and networks on the economy.
Health care
Analysis of social networks is increasingly incorporated into heath care analytics, not only in epidemological studies but also in models of patient communication and education, disease prevention, mental health diagnosis and treatment, and in the study of health care organizations and systems.
Human ecology
Human ecology is an interdisciplinary and transdisciplinary study of the relationship between humans and their natural, social, and built environments. The scientific philosophy of human ecology has a diffuse history with connections to geography, sociology, psychology, anthropology, zoology, and natural ecology.
Language and linguistics
Studies of language and lingustics, particularly evolutionary linguistics, focus on the development of linguistic forms and transfer of changes, sounds or words, from one language system to another through networks of social interaction. Social networks are also important in language shift, as groups of people add and/or abandon languages to their repertoire.
Literary networks
In the study of literary systems, network analysis has been applied by Anheier, Gerhards and Romo, De Nooy, and Senekal, to study various aspects of how literature functions. The basic premise is that polysystem theory, which has been around since the writings of Even-Zohar, can be integrated with network theory and the relationships between different actors in the literary network, e.g. writers, critics, publishers, literary histories, etc., can be mapped using visualization from SNA.
Organizational studies
Research studies of formal or informal organizational relationships, organizational communication, economics, economic sociology, and other resource transfers. Social networks have also been used to examine how organizations interact with each other, characterizing the many informal connections that link executives together, as well as associations and connections between individual employees at different organizations. Intra-organizational networks have been found to affect organizational commitment, organizational identification, interpersonal citizenship behaviour.
Social capital
Social capital is a sociological concept which refers to the value of social relations and the role of cooperation and confidence to achieve positive outcomes. The term refers to the value one can get from their social ties. For example, newly arrived immigrants can make use of their social ties to established migrants to acquire jobs they may otherwise have trouble getting (e.g., because of lack of knowledge of language). Studies show that a positive relationship exists between social capital and the intensity of social network use.
See also
- Collective network
- Complex networks
- Dynamic network analysis
- International Network for Social Network Analysis
- Interpersonal relationship
- Network science
- Network society
- Network theory
- Semiotics of social networking
- Social complexity
- Social group
- Social media
- Social network analysis
- Social networking
- Social relation
- Social web
Further reading
- Wellman, Barry; Berkowitz, S.D. (1988). Social Structures: A Network Approach. Structural Analysis in the Social Sciences. Cambridge University Press. ISBN 0-521-24441-2.
- Scott, John (1991). Social Network Analysis: a handbook. SAGE. ISBN 978-0-7619-6338-7.
- Wasserman, Stanley; Faust, Katherine (1994). Social Network Analysis: Methods and Applications. Structural Analysis in the Social Sciences. Cambridge University Press. ISBN 978-0-521-38269-4.
- Barabási, Albert-László (2003). Linked: How everything is connected to everything else and what it means for business, science, and everyday life. Plum. ISBN 978-0-452-28439-5.
- Freeman, Linton C. (2004). The Development of Social Network Analysis: A Study in the Sociology of Science. Empirical Press. ISBN 1-59457-714-5.
- Barnett, George A. (2011). Encyclopedia of Social Networks. SAGE. ISBN 978-1-4129-7911-5.
- Kadushin, Charles (2012). Understanding Social Networks: Theories, Concepts, and Findings. Oxford University Press. ISBN 978-0-19-537946-4.
- Rainie, Lee and Barry Wellman. 2012. Networked: The New Social Operating System. MIT Press. isbn=978-0-262-01719-0
- E. Estrada, "The Structure of Complex Networks: Theory and Applications", Oxford University Press, 2011, ISBN 978-0-199-59175-6
Peer-reviewed journals
- Social Networks
- Network Science
- Journal of Social Structure
- Journal of Mathematical Sociology
- Social Network Analysis and Mining (SNAM)
Textbooks and educational resources
- Networks, Crowds, and Markets (2010) by D. Easley & J. Kleinberg
- Introduction to Social Networks Methods (2005) by R. Hanneman & M. Riddle
- Social Network Analysis Instructional Web Site by S. Borgatti
Data sets
|Wikimedia Commons has media related to: Social networks|
- Pajek's list of lists of datasets
- UC Irvine Network Data Repository
- Stanford Large Network Dataset Collection
- M.E.J. Newman datasets
- Pajek datasets
- Gephi datasets
- KONECT - Koblenz network collection
- RSiena datasets
- Wasserman, Stanley; Faust, Katherine (1994). "Social Network Analysis in the Social and Behavioral Sciences". Social Network Analysis: Methods and Applications. Cambridge University Press. pp. 1–27. ISBN 9780521387071.
- Scott, W. Richard; Davis, Gerald F. (2003). "Networks In and Around Organizations". Organizations and Organizing. Pearson Prentice Hall. ISBN 0-13-195893-3.
- Freeman, Linton (2004). The Development of Social Network Analysis: A Study in the Sociology of Science. Empirical Press. ISBN 1-59457-714-5.
- Borgatti, Stephen P.; Mehra, Ajay; Brass, Daniel J.; Labianca, Giuseppe (2009). "Network Analysis in the Social Sciences". Science 323 (5916): 892–895. doi:10.1126/science.1165821.
- Easley, David; Kleinberg, Jon (2010). "Overview". Networks, Crowds, and Markets: Reasoning about a Highly Connected World. Cambridge University Press. pp. 1–20. ISBN 978-0-521-19533-1.
- Scott, John P. (2000). Social Network Analysis: A Handbook (2nd edition). Thousand Oaks, CA: Sage Publications.
- Tönnies, Ferdinand (1887). Gemeinschaft und Gesellschaft, Leipzig: Fues's Verlag. (Translated, 1957 by Charles Price Loomis as Community and Society, East Lansing: Michigan State University Press.)
- Durkheim, Emile (1893). De la division du travail social: étude sur l'organisation des sociétés supérieures, Paris: F. Alcan. (Translated, 1964, by Lewis A. Coser as The Division of Labor in Society, New York: Free Press.)
- Simmel, Georg (1908). Soziologie, Leipzig: Duncker & Humblot.
- For a historical overview of the development of social network analysis, see: Carrington, Peter J. & Scott, John (2011). "Introduction". The Sage Handbook of Social Network Analysis. SAGE. p. 1. ISBN 978-1-84787-395-8.
- See also the diagram in Scott, John (2000). Social Network Analysis: A Handbook. SAGE. p. 8. ISBN 978-0-7619-6339-4.
- Malinowski, Bronislaw (1913). The Family Among the Australian Aborigines: A Sociological Study. London: University of London Press.
- Radcliffe-Brown, Alfred Reginald (1930) The social organization of Australian tribes. Sydney, Australia: University of Sydney Oceania monographs, No.1.
- Radcliffe-Brown, A.R. (1940). "On social structure". Journal of the Royal Anthropological Institute, 70, 1-12.
- Lévi-Strauss, Claude (1967). Les structures élémentaires de la parenté. Paris: La Haye, Mouton et Co. (Translated, 1969 by J. H. Bell, J. R. von Sturmer, and R. Needham, 1969, as The Elementary Structures of Kinship, Boston: Beacon Press.)
- Barnes, John (1954). "Class and Committees in a Norwegian Island Parish." Human Relations, (7): 39-58.
- Freeman, Linton C. and Barry Wellman (1995). "A note on the ancestoral Toronto home of social network analysis." Connections, 18(2): 15-19.
- Savage, Mike (2008). "Elizabeth Bott and the formation of modern British sociology." The Sociological Review, 56(4): 579–605.
- Nadel, SF. 1957. The Theory of Social Structure. London: Cohen and West.
- Parsons, Talcott ( 1949). The Structure of Social Action: A Study in Social Theory with Special Reference to a Group of European Writers. New York, NY: The Free Press.
- Parsons, Talcott (1951). The Social System. New York, NY: The Free Press.
- Blau, Peter (1956). Bureaucracy in Modern Society. New York: Random House, Inc.
- Blau, Peter (1960). "A Theory of Social Integration." The American Journal of Sociology, (65)6: 545-556 , (May).
- Blau, Peter (1964). Exchange and Power in Social Life.
- Bernie Hogan. "The Networked Individual: A Profile of Barry Wellman".
- Granovetter, Mark (2007). "Introduction for the French Reader," Sociologica 2: 1–8
- Wellman, Barry (1988). "Structural analysis: From method and metaphor to theory and substance." Pp. 19-61 in B. Wellman and S. D. Berkowitz (eds.) Social Structures: A Network Approach, Cambridge, UK: Cambridge University Press.
- Mullins, Nicholas. Theories and Theory Groups in Contemporary American Sociology. New York: Harper and Row, 1973.
- Tilly, Charles, ed. An Urban World. Boston: Little Brown, 1974.
- Mark Granovetter, "Introduction for the French Reader," Sociologica 2 (2007): 1–8.
- Wellman, Barry. 1988. "Structural Analysis: From Method and Metaphor to Theory and Substance." Pp. 19-61 in Social Structures: A Network Approach, edited by Barry Wellman and S.D. Berkowitz. Cambridge: Cambridge University Press.
- Nagler, Jan, Anna Levina and Marc Timme (2011). "Impact of single links in competitive percolation." Nature Physics, 7: 265-270.
- Newman, Mark, Albert-László Barabási and Duncan J. Watts (2006). The Structure and Dynamics of Networks (Princeton Studies in Complexity). Oxford: Princeton University Press.
- Wellman, Barry (2008). "Review: The development of social network analysis: A study in the sociology of science." Contemporary Sociology, 37: 221-222.
- Faust, Stanley Wasserman; Katherine (1998). Social network analysis : methods and applications (Reprint. ed.). Cambridge [u.a.]: Cambridge Univ. Press. ISBN 0521382696.
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Building code inspector officals have been inconsistant in their reactions to the concerns about indoor air quality.Concerns about energy consumption in the 1970's forced many building code inspectors to severly restrict the amount of outside air used in the building.Indoor air quality concerns have reversed that trend,with some communties now insisting on very large amounts of outside air for wall cavity,and many other building applications.
Building Pressurization- Structural Moisture Outside air intrduced into a building and inside air exhausted to the outside are both forms of ventilation.Properly using and balancing each type of air is essential.Stack effect and wind influence a building's internal pressure.Stack effect,caused by buoyant warm air rising to the top of the building surronded by cold air,increases the pressure at the top of the building envelope.The bottom of the building has less pressure than the top as the warm air inside rises from it.The building expells warm air at its top(because high pressure there)while inhaling cold outside air at its lower pressure.A properly operating ventilation system should be setup to slightly pressurize the interior space.This is not possible under all conditions.The goal,however,of a property operating ventilation system is to provide building pressure under most conditions.There are many benefits to slight postive pressure.Reduction in air infiltration into the home like uncontrolled drafts and dust intake.Positive pressure is attained by properly balancing the exhaust and intake air supplies to the building.Ventilation air is required throughout the year,including when tempertures are well below freezing.
Humidity Control-Proper regulation of humidity is of increasing concern to control condensation in wall cavitys,and to assure good air quality.Air that is to moist or dry can lead to serious problems such as mold growth.Living in the south east where humidity levels are sometimes in the 80 to 90% range on certain days.Damage to the building structure and occupants illness can be caused by improper humidity control.Indoor air quality-the existence of bacteria,mold,and other micoorganisms in the air-is often related to humidity control.What relative humidity should I have in my home? Seems like a simple enough question. However, the answer can sometimes be difficult to understand.
Elevated relative humidity at a surface – 70 percent or higher - can lead to problems with mold, corrosion, decay and other moisture related deterioration. When relative humidity reaches 100 percent, condensation can occur on surfaces leading to a whole host of additional problems. An elevated relative humidity in carpet and within fabrics can lead to dust mite infestation and mildew (mildew is mold growing on fabrics).
Low relative humidity can lead to discomfort, shrinkage of wood floors and wood furniture, cracking of paint on wood trim and static electricity discharges.
The key is not to be too low and not to be too high. High enough to be comfortable, but low enough to avoid moisture problems associated with mold, corrosion, decay, and condensation.
Unfortunately, determining the correct range depends on where the home is located (climate), how the home is constructed (the thermal resistance of surfaces determines surface temperatures), the time of year (the month or season determines surface temperatures), and the sensitivity of the occupants.
Limits to Relative Humidity—Comfort and Health Aspects
How low can you go? Comfort wise at least, the 2001 ASHRAE Fundamentals (8.12) tells us that at dew point temperatures of less than 32 degrees F, complaints of dry nose, throat, eyes, and skin occur. A dew point of 32 degrees converts to a relative humidity of 25 percent at 68 degrees.
How high can you go? Again, using comfort as the criteria, the 2001 ASHRAE Fundamentals (8.12) tells us that a relative humidity of 60 percent should not be exceeded.
This is consistent with ASHRAE Standard 62-2001 Ventilation for Acceptable Indoor Air Quality, which recommends that the lower boundary of the relative humidity range be limited to 25 percent and the upper boundary of the relative humidity range be limited to 60 percent.
Now, it is important to consider the ASHRAE definition of comfort: “combinations of indoor space environment and personal factors that will produce thermal environmental conditions acceptable to 80 percent or more of the occupants within a space.” Remember that you can’t please all of the people all of the time.
The ranges cited above do not consider health, except indirectly. Some people love to live in desert climates, and some people love to live in the tropics. The upper limits from a health perspective are indirectly derived from a desire to control the growth of mold, bacteria, and other disease vectors. Similarly, for the lower limits, although the lower limits tend to be arguably “healthier” from a disease vector perspective. Dry conditions do not favor mold, most bacteria, and other disease vectors.
However, some have argued that dry conditions dry out the mucus linings of the respiratory system and therefore make it more difficult for the body to fight off invaders. The other side of the argument is that there are fewer invaders to worry about.
As can be expected, individual sensitivities and susceptibilities vary greatly, and it is typically very difficult to generalize with respect to relative humidity and health. Having said it is difficult to generalize, we will do so anyway. Keeping relative humidity in the 25 percent to 60 percent range tends to minimize most health issues – although opinions vary greatly.
Incorrect recommendations in the popular press often lead occupants and homeowners to over humidify homes during the winter. The range of 40 percent to 60 percent relative humidity is commonly incorrectly recommended for health and comfort reasons. As we will see, there is a big difference between 25 percent as a lower limit rather than 40 percent – particularly in very cold and cold climates.
To complicate things further, most people are not capable of sensing relative humidity fluctuations within the range of 25 percent to 60 percent. If the relative humidity drops below 25 percent, most people can sense it. Similarly, if the relative humidity rises above 60 percent most people can sense it. In the range of 25 percent to 60 percent the majority of people cannot sense any difference. The range of 25 percent to 60 percent is typically defined as the comfort range for this reason. This is very different than people sensing temperature variations. Most people can sense a difference in temperature within a range of 1 to 2 degrees. Less—below 0.1 degrees—if they are married (just kidding).
Comfort is of course different than health. When relative humidity drops below 25 percent there have been some reports in the medical literature of eye irritation in office workers using computers. Breathing difficulties have been reported in some individuals when relative humidities drop below 15 percent due to the mucus linings of the respiratorysystem
desiccating. However, there is no medical consensus in this regard.
Determining the Humidity Limits—
Many people believe that 25 percent relative humidity as a lower level is still too high. The debate breaks predictably into several camps with the engineers (the aircraft people being the most vocal) arguing for no lower limit for health and only a discussion on comfort. Whereas the lung researchers and some MD’s argue that until there is definitive research, why not keep the level high from a prudent avoidance perspective. This of course terrifies the microbiologists and mold researchers since higher lower limits clearly lead to mold growth in buildings and are associated with microbial contamination in typical residential humidifiers.
So on the lower limit there is no real consensus, but only a current compromise recommendation. It is pretty clear that the lower limit will not go up. The only question is how low it will end up. At present, 25 percent relative humidity is the current compromise recommendation within ASHRAE.
On the upper end, there is an emerging consensus. Interior relative humidity should be maintained so that a 70 percent relative humidity at a building surface is avoided in order to control mold growth and should never rise above 60 percent in any event.
Relative Humidity, Surface Humidity, and Condensation
Consensus among microbiologists gives the critical relative humidity for adverse biological activity to occur on building envelope surfaces to be 70 percent. Where a relative humidity above 70 percent occurs at surfaces, mold growth, dust mite growth, decay, corrosion, etc. can occur. Therefore, conditions should be maintained within a building such that the critical 70 (or higher) percent relative humidity at a building envelope surface does not occur. Due to climate differences, interior conditions which must be maintained to avoid the critical relative humidity at a surface vary from region to region and time of year. They also vary based on the thermal resistance of the building envelope.
This means in winter months in cold climates interior relative humidity should be kept as low as possible but within the comfort and health range (i.e. above 25 percent if you believe ASHRAE Standard 62-2001).
In the summer months it means that interior relative humidity should never exceed 60 percent for both comfort and health reasons.
There is a fundamental difference between relative humidity measured in the middle of a conditioned space, and the relative humidity found at surfaces due to the significant difference in temperature typically found between surfaces and the air in the middle of a conditioned space.
For a given sample of air containing water, relative humidity goes up as the temperature goes down. If the air in the middle of a room is 70 degrees at a relative humidity of 40 percent, any surface below 45 degrees will be able to condense water. Any surface below 54 degrees will have air adjacent it at a relative humidity of 70 percent – the mold limit.
Whereas when air in the middle of the room is 70 degrees at a relative humidity of 25 percent, the temperature of a condensing surface drops to 32 degrees from 45 degrees. And a surface with a relative humidity adjacent to it of 70 percent drops to 40 degrees from 54 degrees.
In other words, for condensation to occur with air at 70 degrees and a relative humidity of 25 percent, surfaces need to be colder than 32 degrees. For mold to grow, surfaces need to be colder than 40 degrees. Of course, in a nice and happy coincidence, mold does not like to grow at surfaces below 40 degrees, but will happily grow at 54 degrees. What does this tell us? Well, if surfaces are likely to be cold – say like in the winter - you are better off having a lower relative humidity.
Where relative humidities near surfaces are maintained below 70 percent, mold and other biological growth can be controlled. Since relative humidities are dependant on both temperature and vapor pressure, mold control is dependant on controlling both the temperature and vapor pressure near surfaces.
Surface Humidity and Building Assemblies: Applications in Heating Climates
In heating climates, mold growth on interior surfaces occurs during the heating season because the interior surfaces of exterior walls are cool from heat loss and because moisture levels within the conditioned space are too high. Mold growth control is facilitated by preventing the interior surfaces of exterior wall and other building assemblies from becoming too cold and by limiting interior moisture levels. The key is to prevent relative humidities adjacent surfaces from rising above 70 percent. The thermal resistance of the building envelope and the local climate determine the interior surface temperatures of exterior walls and other building assemblies. Controlled ventilation and source control limit the interior moisture levels.
Experience has shown, that where interior moisture levels in very cold climates during the heating season are limited to the 25 percent relative humidity at 70 degrees, relative humidities adjacent to the interior surfaces of exterior walls (of typical code minimum thermal resistance) fall below 70 percent and mold growth is controlled. The colder the climate (for the thermal resistance of any given building envelope) the lower the interior relative humidity necessary to prevent 70 percent relative humidities occurring adjacent interior surfaces of exterior walls. Building enclosures of similar thermal resistance interior moisture levels during the heating season. A 25 percent interior relative humidity at 70 degrees would be appropriate for Minneapolis. Whereas interior relative humidities up to 35 percent at 70 degrees would be appropriate for Cincinnati – which is located in a cold climate rather than a very cold
climate . Correspondingly, the higher the desired interior relative humidity, the higher the
thermal resistance necessary to control relative humidities adjacent to interior surfaces.
In a mixed climate, during the heating season, interior moisture levels should be limited to 45 percent relative humidity at 70 degrees. This limits the relative humidity adjacent to the interior surface of exterior walls to below 70 percent for the typical thermal resistance found in most building assemblies in this climate zones.
In cooling climates, interior mold growth also occurs because interior surfaces are typically cold and then exposed to moisture levels that are too high. The cold surfaces in cooling climates arise from the air conditioning of enclosures. When exterior hot air is cooled, its relative humidity increases. If the exterior hot air is also humid, cooling this air will typically raise its relative humidity above the point at which mold growth can occur (70 percent).
Where air conditioned "cold" air is supplied to a room, and this air can be "blown" against an interior surface due to poor diffuser design, diffuser location, or diffuser performance, creating cold spots on the interior gypsum board surfaces. Although this cold air is typically dehumidified before it is supplied to the conditioned space, it can create a mold problem on room surfaces as a result of high levels of airborne moisture within the room contacting the cooled surface. This typically leads to a rise in relative humidity near the surface and a corresponding mold problem.
If exterior humid air comes in contact with the interstitial cavity side of cooled interior gypsum board mold and other biological growth can occur. Cooling this exterior hot, humid air by air conditioning or contact with cool surfaces will raise its relative humidity above 70 percent. When nutrients are present mold and other growth occurs. This is exacerbated with the use of impermeable wall coverings such as vinyl wallpaper that can trap moisture between the interior finish and the gypsum board. When these interior finishes are coupled with cold spots (from poor diffuser placement and/or overcooling) and exterior moisture, mold and other growth can occur.
Accordingly, one of the most practical solutions in controlling mold and other biological growth in cooling climates is the prevention of hot, humid exterior air, or other forms of moisture transport, from contacting the interior cold (air conditioned) gypsum board surfaces (controlling the vapor pressure at the surface). This is most commonly facilitated by maintaining the conditioned space at a positive air pressure to the exterior and the installation of an exterior vapor diffusion retarder. Pressurization of building assemblies
is expedited by airtight construction.
Interior moisture levels within the conditioned space should also be limited to 60 percent relative humidity at 75 degrees by dehumidification and source control to prevent mold growth on the interior surfaces within the conditioned space.
Many people are concerned about wood floors and wood furniture being damaged if humidifiers are not installed. More often than not, people tend to over humidify their homes in an attempt to protect their wood floors and wood furniture. They need not do so if relative humidities are maintained in the range of 25 to 60 percent between winter and summer.
Let us examine the effect of varying humidity inside of a home between a low of 25 percent and a high of 60 percent on wood. Wood moisture content changes directly with exposure to varying relative humidity.
The relationship is extremely well understood by generations of wood workers and furniture makers .The moisture content of wood will vary from 5 percent moisture content by
weight at 25 percent relative humidity to 11 percent moisture content by weight at 60 percent relative humidity. This results in a maximum change in dimension of approximately 2 percent tangential to the grain. If the wood in question is oak, and the board is 4 inches wide, the
maximum movement is 0.08 inches.
If we have a wood floor installed with 4 inch wide wood boards initially conditioned to the mid range of expected moisture content, i.e. 8 percent moisture content by weight, the range in movement is plus and minus 0.04 inches or approximately the thickness of a credit card. This is not an aesthetically displeasing or unacceptable range of movement. Of course if the wood is not initially conditioned to the mid range of the expected moisture content, then the movement can be two credit card thicknesses.
The amount of bound water in wood is determined by the relative humidity (RH) of the
surrounding atmosphere; the amount of bound water changes (albeit slowly) as the
relative humidity changes. The moisture content of wood, when a balance is established
at a given relative humidity, is its equilibrium moisture content (EMC). The solid line
represents the curve for white spruce, a typical species with fiber saturation point (FSP)
around 30% EMC. For species with a high extractive content, such as mahogany, FSP is
around 24%, and for those with low extractive content, such as birch, FSP may be as high
as 35%. Although a precise curve cannot be drawn for each species, most will fall within
the color band. As a Certified IAQ Inspector I will inspect the check relative humidity levels.Serving atlanta metro and south atlanta counties.Member international code council,and member in good standing with the National Association of Certified Home Inspectors.Ventilation is the key to a health home.
Air conditioning systems that don't control inside humidity leave the space feeling cool and damp.Water vapor from the ouside infiltates the interior when the refrigeration compressor isn't on.Moisture enters with ventilattion air and leakage around the building's perimeter.If the air condition does run almost constanly,interior relative humidity rises.Cooling system that are to small for their application can also benefit from the use of a dehumidifier.Dehumidifiers cannot,however,compensate for a poorly constructed building.A badly or incorrectly installed vapor barrier can admit much more water than must dehumidification system can remove.These areas are very important areas needing inspected by a certified home inspector.
Air quality in buildings with little outdoor air supply quickly deteriorates when the interior space is occupied.Carbon dioxide gas exhaled by people accumulates,leading to discomfort and drowsiness.Poor ventilation does not expel the carbon dioxide or dilute it,and its concentration rises throughout the day.Some quality recovery usally takes places at night while the space is vacant.Other ventilation problems related problems include the accumaltion of gases and dust generated from sources within the home.Many building materials give off chemicals,especially when new.Formaldehyde is commonly "exaled by carpeting,upholstered furniture,and many wood fiber products.This gas is very irritating to the eyes,nose,and throat..People allerigic to it can be sticken with severe reactions that might require hospitalation.
Dust and dirt can spread to with poorly designed or serviced air handling equipment.Improper filtration allows irritating from one part of the homes interior other parts of the home..
Sources of Attic Mold: Roof leaks or, alternatively, high levels of attic moisture due to a combination of inadequate attic (soffit intake and ridge outlet) ventilation combine with building moisture sources (such as a chronic or even a single-event wet basement, plumbing leaks, or a leaky roof from roof failure or from ice dams) are likley to cause excessive moisture or actual wet conditions in an attic. High attic moisture levels or actual wet attic conditions invite extensive mold growth.
Visible mold may appear on wood surfaces in an attic such as on rafters or roof sheathing. Hidden mold may be present and may be even more of a problem if it forms in insulation or in the ducts and air handler of an air conditioning or heating/air conditioning system. Typical building air convection currents tend to move air up and out from lower to upper building levels, so one would not think that much mold would move down from an attic into the living area. But important exceptions to this can quickly move problem mold from an attic into a living area. Conditions moving mold downwards from an attic include
*Mold growth in HVAC ducts or air handlers found in an attic:*
Mold on any attic surface or in attic insulation if it is a species producing airborne spores and if the building uses a whole house ventilating fan, especially if there is inadequate exit venting for the fan operation. This condition pressurizes the attic and moves mold down through various openings into the floors below.
Mold on building surfaces in an attic or attic knee wall space which opens onto or has a knee wall common with an upper floor living space such as a bedroom.
Building Exteriors Leaks and Mold No mold cleanup project will be successful unless you correct the conditions that caused mold growth in the first place. An expert inspection and report should find and suggest remedies for site and building exterior conditions that produce mold or for building areas that serve as a mold reservoir or as amplifiers for allergens, mold, mildew, excessive pollen or pet dander, The basic steps: find all unwanted moisture sources, correct appropriate building, site, landscaping, & construction details. 90% of the wet basements and crawl spaces I see are caused by bad or missing roof gutters and downspouts.
Cleaning mold from smooth wood surfaces
Perfectly adequate cleaning may be accomplished by wiping or (where feasible) power-washing or media blasting. Where wiping a moldy surface, take care not to spread moldy debris from a moldy surface onto a previously uncontaminated surface by making the mistake of re-using the same moldy rag over and over on all surfaces. Professionals use "steri-wiping" which takes care to avoid spreading moldy debris by always folding and using a clean side of the wipe when moving to a new spot.
Where the framing lumber is indoors or otherwise in a location where water spillage is a concern, wipe the areas of heaviest mold to remove any loose mold from the surface of the lumber. Unless professional area-containment has been set up (barriers, negative air), do not use violent cleaning methods such as power-washing or sandblasting indoors, as you will spread moldy debris throughout the building and you'll increase the ultimate project cleanup cost. (Where the framing lumber is outdoors where water spillage and the creation of aerosolized mold spores is not an issue, pressure wash the infected lumber to remove surface mold.
"Cleaning" in this case can be simply wiping with a sponge wet with water or detergent. See my warnings below about using bleach. The object of cleaning is to remove most of the loose moldy particles. The object (except in medical facilities) is not to produce a particle-free sterile surface. However beware of cross-contamination. Wetting a rag and wiping a very moldy surface off is fine but if you then use the same dirty rag to wipe another fairly clean surface you may be in fact spreading moldy debris around. A professional uses sterile wipes and folds to a clean side of the wipe for each wiping stroke. For a small homeowner non-critical project this may be overkill but think about and avoid spreading moldy debris by your cleaning procedure.
"We Look Beyond the Basics" |
Globalization and its Effects on the Transnational Movement to End Human Trafficking
Globalization and its Effects on the Transnational Movements to End Human Trafficking
Globalization is known as a phenomenon occurring all around us, making positive contributions to our lives, as well as negative. In addition to fueling the technology revolution, globalization has made us closer to achieving such goals as free trade, deregulation, and economic flexibility. At the same time, however, this time-space compressor is seen as a monster to millions of people around the world; it is a principle cause of the growing inequality gap, of political domination, poverty and debt, as well as, oppression and structural violence.1 With that being said, it is clear that globalization does not, in fact, make everyone’s grass greener. However, particularly the main benefiters of globalization such as the Western nations, many people do not realize the dramatic negative effects of this phenomenon. Therefore, one of the ways to better your understanding of the positive and negative effects of globalization is by expanding their knowledge of a transnational social movement.
Thousands of movements exist today, each involving intersecting layers of networks within hundreds of nations across the globe. Various movements exist to fight the many forms of oppression and violence that have become more prominent since globalization took a more dominant role in shaping our world. A transnational crime that is being affected in an incredibly negative way by globalization is that of modern-day slavery, better known as human trafficking. For instance, worldwide, 27 million people are victims of human trafficking and, with its annual revenue between 5 and 7 billion dollars, it is clear that trafficking is an enormous problem rendering global attention.2 Since globalization plays such a primary role in the transnational anti-human trafficking movement, as well as in all transnational social movements, in this paper I will first discuss the relationship between globalization and transnational social movements in general, discussing some of the positive and negative effects. Then, I will elaborate furth er on human trafficking (specifically sex-trafficking) as a major global problem, as well as, delve into how the anti-trafficking movements are dominantly effected by globalization. Within the context of that, I will highlight the problem of sex-trafficking in India (a country where it is most prevalent), and finish by focusing on the Coalition Against Trafficking in Women (CATW), the first non-governmental organization to fight human trafficking (especially sex trafficking in women and girls).3 Inspired by the CATW, while addressing human trafficking as a whole in some aspects, this paper will focus explicitly on the fights to end the sex-trafficking of women and girls.
Similar to transnational movements, Valentine M. Moghadam stated in her book, Globalization and Social Movements, that “globalization can be thought of as something that “transcends nation state-boundaries; a multi-faceted process of social change with economic, political, and cultural dimensions, creating new forms of inequality, competition, and transnational forms of organizing and mobilizing.”4 A transnational social movement is something that unifies people, in three or more countries, to work together on an international issue that they feel passionately enough about to want to come together and fight.5 Today, some of the transnational social movements with the largest number of supporters are: feminist movements, gay/lesbian rights movements, and Islamic movements. With that being said, the ultimate goal of any social movement is to delegitimize their enemy’s legal stance, eventually leading to their demise. However, this becomes increasingly more difficult when the primary force that a movement is fighting is something as eternal as globalization.
Nonetheless, globalization has its positive and negative effects to all social movements. First, I will discuss one of the positive effects that globalization has had on the world, specifically in terms of aiding transnational social movements. Technology, ironically also a negative effect, has drastically expanded the range of transnational social movements with things such as the internet and the cell phone. It is through the use of technology that movements and organizations (especially transnational) are able to recruit a lot of their supporters. With the internet comes instant communication and advertisement, making contact between existing members of a group, as well as future members, much easier. In addition to facilitating contact and recruitment, the internet serves as a great avenue for raising global awareness. A lot of transnational organizations may not have an enormous formal following but that does not mean that they have been unsuccessful in their recruitment endeavors. Someone does not have to be an official member of a transnational organization or social movement to be able to contribute to spreading global awareness about a particular injustice. In fact, some of the most successful advertising is through word of mouth, something that all of us are capable of contributing to.
Another aspect about transnational social movements that the global technology revolution has contributed to is the expansion of their goals and audience. According to Kevin Bales, author of Understanding Global Slavery, “the average non-state organization, before globalization, aimed only to alter state policies, usually reaching national borders at the very most.”6 Now, social movements are not forcibly restricted (except if by their own personal choice) to state boundaries. As a result, their target audiences have expanded exponentially, originally being primarily local and becoming global. Just by basing their goals on philosophies relatable to all audiences and making informational resources available, people can find out information faster, spread awareness quicker and easier, and even donate to a cause in a matter of seconds. This creates a lot higher of a success rate for organizations with less required work, such as campaigning, protesting, and personally asking people to donate. All of these options would not be presented to transnational social movements without the fuel of technology, which would not be nearly as prominent and available in our world without globalization.
Now taking a look at the other side of globalization, we will explore some of its negative effects, including overpopulation and the growing inequality gap in addition to discussing how technology is also a negative effect of globalization.7 Beginning with overpopulation, globalization has played a huge role in not just overpopulation and migration but in minimizing the resources available to third world impoverished nations. For instance, after the population explosion of WWII raised the global population from 2 to 6 billion, problems regarding global population began to exponentially advance.8 When large numbers of impoverished people inhabit a nation, it is often that they find themselves under the control of a corrupt government. Thus, when the people that are supposed to be protecting them, particularly law enforcement, are in conspiracy with the enemy (such as transnational criminal organizations), citizens are not able to protect themselves against things such as enslavement, trafficking, discrimination, and several other forms of oppression and violence. This makes transnational social movements hopes of success increasingly more implausible. The more the population grows, the more people are likely to be subject to transnational violence and crime, creating a rewinding effect to their already small progress.
Since globalization is not rendering benefits to all nations equally, with Western and other industrialized nations receiving the most advantages, globalization becomes much more negative than positive. With that in mind, the second negative aspect of globalization is the growing economic inequality gap; the rich seem to be getting richer while the poor are getting poorer. This effect of globalization is viewed by most as a direct result of the unequal distribution of its benefits. With such inequality comes an increased amount of poverty, crime, structural violence (oppression), civil wars, governmental corruption, illiteracy, and fatality. To first expand on the effects of poverty, when the poor population of a nation grows, more people are forced, particularly women, to travel internationally for work in order to support their families. In turn, this makes them more vulnerable to being a victim of transnational crime. Also, when in times of immense need, “some parents will sell their children, not just for money but also in hope that the children will be escaping a situation of poverty and will move to a place offering more opportunities.”9 Unfortunately, whether through good intentions or not, those children often grow up without any education or social guidance. Without those, they become increasingly more likely to either become an instigator or a victim of violence (such as ethnic violence, trafficking, invasions, civil wars, etc.). This directly effects transnational social movements, especially those fighting forms of oppression and violence, because it is drastically increasing the amount of victims of such atrocities as trafficking and exploitation.
To further expand on the effects of governmental corruption caused by globalization, when a government is corrupt, the citizens of a nation have no one to trust. When a country’s citizens cannot even trust their own government, crime and structural violence become even more likely to occur. As such rates increase, civil wars, ethnic violence, or invasions often result, causing innocent civilians to lose their home, risk sexual assault or exploitation, lose their children (often to the recruitment into the army as child soldiers), or even die. Therefore, through these negative effects, globalization is breeding a domino effect, beginning with overpopulation, leading to poverty, influencing an increase in crime and governmental corruption, which in turn lead to structural violence, civil wars, and an increase in fatalities. Since they are placing such a great amount of people at risk for their human rights to be intrinsically violated, it is clear that none of these effects pose positive aspects for transnational social movements.
The last negative effect of globalization to be discussed is technology. Although there are also many positive attributes of technology, which were previously discussed, there are even more negative consequences associated with globalization and the technology revolution. For instance, according to Kevin Bales, “the technology and transformationalism that [inspires] the key globalized industries [permits] other types of [negative] activities, [such as] money laundering and trafficking, to assume a global scale.”10 On that note, technology has served as a primary vehicle in making local crimes national, national crimes transnational, and transnational crimes global phenomena, due to the window of opportunity it presents for instant and mass communication. Although such aspects can also be positive, with the ability to communicate an infinite amount of information, whether true or false, to any amount of people in a matter of seconds creates a large, and almost alluring, invitation for new forms of manipulation. Especially since the internet is public domain for anyone to post whatever they want, millions of people fall under the traps of con-artists and criminals, ending up as victims of fraud or future victims of trafficking, money laundering, or sexual or physical exploitation in a matter of one click. Not only does the internet serve as an opportunity for perpetrators to communicate with their future victims, but it also permits members within transnational crime organizations, such as the Taliban, to communicate with ease. Basically, although it has negative and positive effects, technology is known as something that has been fueled by globalization and, in turn, serves as a provider to a great deal of the manipulation and transnational crime in our world today.
Now I will begin expanding on why the global issue of human trafficking is so difficult to fight, by first emphasizing that, although I will be focusing on just one, there are many different types of trafficking. In my opinion, human trafficking (in all of its forms) is a prime example of how history can repeat itself or, in some cases, problems can continue to exist but are recognized as a different crime than in the past. For instance, human trafficking is often referred to as modern-day slavery. However, trafficking is not seen as identical to slavery, since it involves the possibility of many different things Before the Anti-Trafficking Protocol in 2000, “there was no international definition, making there no physical basis for research on the problem.”11 Now, a victim of human trafficking is commonly understood as someone who is taken against their will and forced to do certain things (such as prostitution or drug laundering) and be subject to things that often compromise the following plus more: their right to free will, right to freedom from slavery, torture, and other inhumane or degrading treatment, the right to freedom from discrimination, the right to human dignity, and the right to work in just conditions. However, that still leaves an incredible amount of room for ambiguity. With that being said, one of the biggest problems with being able to prevent human trafficking is that since there are so many avenues by which someone can be a victim of trafficking. Hence, people find it difficult to pinpoint the crime into specifics. With such a narrow focus placed on sex trafficking, sometimes the word “trafficking” is used incorrectly, directly referring to the transportation and sexual exploitation of victims (aka sex trafficking).12 Although it is definitely one of the most common types of trafficking, sex-trafficking should not be thought of as the only form of trafficking that is prevalent today. To help expand the spectrum, some forms of trafficking that are also current transnational issues are: trafficking for the purpose of agricultural work or other forced labor, drug trafficking, and weapon trafficking; each of which have millions of victims annually. Since attempting to tackle anti-tracking as a whole would be far too difficult to discuss all at once, although all forms of trafficking deserve global attention, this paper will specifically focus on sex-trafficking (spotlighting women and girls as victims) and the transnational social movements to fight it.
Now that the basic definitions and understandings of the various types of trafficking have been successfully established, I will continue discussing in more detail the severity of sex-trafficking, beginning with more aspects of sex-trafficking and why fighting it is so difficult, following with an expansion the anti-trafficking social movement in relation to globalization, spotlighting the issue as it takes place in India, and finishing with an exploration of the specific transnational social organization: Coalition Against Trafficking in Women (CATW). On that note, another one of the most detrimental aspects of the sex-trafficking phenomenon is that it is very difficult to determine how severe it actually is, so by default, we assume that it is drastically underestimated. To expand on this a little, although some people believe that NGOs have inflated statistics to influence anti-trafficking policies, based on current research on other forms of sexual assault and exploitation (such as rape), it is highly unlikely that the number of sex-trafficking victims is overestimated. One of the main reasons why it is so difficult to determine its severity is that it is extremely challenging to gain any information from former or current victims or traffickers. Furthermore, to create a secretive nature similar to other sexually-evasive crimes, traffickers employ certain tactics on their victims to prevent them from escaping or being found and reporting the crime. Sex-traffickers often threaten their victims with violence, deportation, or murder; force them to take an excessive amount of drugs (almost always causing addiction and often to the point of an overdose); or they just make them feel so ashamed about their body and who they’ve become, beating down their self-esteem, that the victims reach the point where they no longer desire to escape. As if that is not enough, traffickers will also take it upon themselves to do anything and everything to keep their victims hidden; since sex-trafficking is a crime in which “the victim is also the moneymaking ‘product,’ like that of a bag of cocaine that a drug trafficker would keep hidden at all costs,” sex-trafficking victims will be exhausted until they are deemed worthless and then disposed of.13 A lack of initiative on the part of the people can also be blamed for fueling the fire of sex-trafficking. Since a great deal of ambivalence surrounds the ideas behind prostitutes and undocumented immigrants (victims of sex-trafficking in disguise), people are slow to take action to protect them.14 With all of these aspects placing a blanket of disguise of sex-trafficking victims and statistics, the transnational social movements to end human trafficking live on with minor noticeable success.
Although this is an issue that takes place worldwide, I wanted to highlight the issue of sex-trafficking in India, as it is one of the world’s most prevalent countries for trafficking and slavery. Starting with some statistics, on a scale of 1-4, with 1 being low and 4 being very high, India ranks 4 in incidences of slavery, 3 as a destination country for trafficking, and 4 as a transit country for trafficking. To better place things into perspective, “the number of slaves in India is estimated to be between 18 and 22 million,” more than any other country by nearly 10 million.15 Although this includes all forms of slavery and trafficking, not just sex-trafficking, but it is without a doubt that with those types of numbers millions of victims are being subject to sexual exploitation as well. Since India sufferers from a lot of globalization’s negative effects, the global consequences that were discussed earlier start to become more apparent.
For instance, India suffers immensely from overpopulation, governmental corruption, and poverty. As an example, it was reported on the website of the transnational social movement against sex-trafficking, CATW (to be discussed later on), that “top politicians and police officials in Bombay are in league with the mafia who control the sex industry, exchanging protection for cash payoffs. many politicians view prostitutes as an expendable commodity.”16 Although there have been minor signs of progress since, the UN Convention of the Suppression of the Traffic in Persons and the Exploitation and the Prostitution of Others was in fact signed by India.17 In terms of the future, efforts to rehabilitate victims and prevent future victims of sex-trafficking are squandered by the lack of support from the Indian citizens and their government. Like all other places where this phenomenon is most prevalent, India needs to raise public awareness on the issue of sex-trafficking and accepting that its negative effects not only influence the Indian population but the world as a whole.
Now that we have seen how devastating sex-trafficking can be in just one country, it should be clear that the issue is an even larger problem when looked at from a global perspective. Thus, I am going to delve into the details of the CATW, the first non-governmental organization dedicated to promoting women’s rights by internationally combating sexual exploitation in all of its forms.18 Within exploring the organization as a transnational movement, I am going to discuss its direct affiliations with globalization. Before evaluating the Coalition Against Trafficking in Women, it is important that two other examples of effort set forth to combat trafficking are briefly mentioned; the first being the Anti-Slavery Society and the second (which was already previously mentioned) being the Anti-Trafficking Protocol. It is critical to recognize that since the Anti-Slavery Society was before globalization became a global concern, the organization was primarily located within the state political system, never reaching international borders.19 Clearly this has changed, since several if not most anti-trafficking organizations today are transnational. However, the goals of the movement have not changed since pre-globalization, just the ways in which the goals are sought to be achieved. For example, the goal of the Anti-Slavery Society was to “bring about public redefinition of slavery as a moral issue, not locally, but globally,” which is the goal of anti-trafficking organizations today, such as the CATW. The only difference is that with the internet and other technological avenues, public awareness is spread easier and faster. Unfortunately, the issue has not improved much since the times of the Anti-Slavery Society because although technology makes for easier communication, as previously discussed, it also makes for easier manipulation and, therefore, more victims. In terms of the Anti-Trafficking Protocol of 2000, besides establishing a more formal definition of trafficking, it also helped the world to recognize that “trafficking [is] a transnational crime requiring a transnational solution, and that globalization and new technologies [have] created new opportunities for criminal organizations.”20
With inspiration from the Anti-Slavery Society and the Anti-Trafficking Protocol, the Coalition Against Trafficking in Women (or CATW), founded in 1988, dominates the international scene today as the most “vocal and best known organization” for advocating for women’s rights, particularly in the form of anti-sex-trafficking campaigns.21 As a leading organization, the CATW has successfully set up worldwide networking against sex-trafficking and prostitution through their major world regions in Africa, Asia, Australia, Europe, Latin America, and North America. They also have national coalitions in over fifteen countries, including the Philippines, Thailand, Venezuela, the United States, Canada, and France. For an organization that is primarily anti-globalization (or at least against the negative effects of globalization), they use the internet in a lot of their work. Aspects such as communication between members, spreading awareness, and recruiting supporters, are all immensely aided by the internet and the use of cell phones. However, the CATW has alternative forms of such communication, such as through their presence at events similar to the World Social Forum, as well as, through word of mouth.22 In addition, CATW works with national and international policy makers, women’s rights and human rights advocates, national congress, and regional and UN committees and commissions.23 ? At a more local level, the CATW strives to educate young boys and girls in schools throughout various parts of the world by training teachers, police, and governmental officials about the harms of sexual exploitation and ways to resist and combat it.24 Thus, with its goals, methods of discourse, recruitment, and techniques, it is clear that the CATW is a very well established transnational social movement; one that seems to be fighting globalization but, at the same time, could not be nearly as successful without it. Hence, this movement can definitely be thought of as an example of globalization from below; in other words, “a reaction to the effects of neoliberal globalization and an exemplar of the transnational and collective action” used to combat it in the form of NGOs, social movements, and civil society organizations.25
In conclusion, it should now be without question that trafficking (particularly sex trafficking) is an immense and growing problem in the globalized world; one that has been helped, but mostly hindered, by the inevitable effects of globalization. Due to the unequal distribution of the its benefits, overpopulation, poverty and crime caused by globalization lead to governmental corruption, causing oppression, structural violence, famine, civil wars, and in particular, atrocious violations of human rights such as sex-trafficking. Before closing, I want to address the fact that many people, particularly citizens of wealthy Western nations, argue that trafficking is only a phenomenon in impoverished and third world countries and that it does not, in fact, place us at risk. On that note, it needs to be known that “since organized criminal groups are alleged to dominate trafficking, it is seen as a [serious] security threat to Western countries.”26 Furthermore, the US is also a very popular destination country for trafficking; just because it is not as severe as it is in countries such as India and Thailand does not mean that it does not exist. Therefore, as citizens of what is arguably the most powerful nation in the world, we cannot ignore our opportunity to make a difference in the fight to end sexual exploitation and other forms of trafficking. As seen in other arenas besides trafficking, globalization can transform local issues into global phenomena very rapidly, so although trafficking may not be as prevalent in the US as it is in developing countries, there is no way of knowing when that might change. Globalization has been known to have no boundaries, making the actions of each country have direct negative effects on the rest of the world. With that being said, I hope to have inspired people to become involved in an anti-trafficking organization, such as the CATW, or to somehow join the global movement to end sex-trafficking and other forms of trafficking, not only in third world countries but worldwide. As the great Gandhi once said, “be the change you wish to see in the world.”27
Bales, Kevin. Understanding Global Slavery: a Reader. Berkeley [u.a.: Univ. of California, . Print. 2005.
“Coalition Against Human Trafficking in Women.” Human Trafficking-Trafficking of Humans-Coalition Against Trafficking of Women. Web. 18 May 2010. .
Human Trafficking Statistics. Polaris Project. Web. 18 May 2010. .
Kyle, David, and Rey Koslowski. Global Human Smuggling: Comparative Perspectives. Baltimore: Johns Hopkins UP, 2001. Print.
Moghadam, Valentine M. Globalization and Social Movements: Islamism, Feminism, and Global Justice Movement. Lanham: Rowman & Littlefield, 2009. Print.?Sen, Amartya. “How to Judge Globalism.” The American Prospect. Web. 18 May 2010. .
Kelly Dona (2011)
Peace Studies & French
Along with being a part of the Honors Program, Kelly Dona will be graduating this May 2011with a double major in French and Peace Studies. After graduation, she hopes to volunteer with the Peace Corps and then pursue a Master’s degree in Peace and Conflict Studies. |
Affected System: Multi system
|Synonyms and Related Keywords|
Cadmium Toxicity, Cadmium Intoxication
Cadmium has no constructive purpose in the human body. It and its compounds are extremely toxic even in low concentrations, and will bioaccumulate in organisms and ecosystems.
The itai-itai disease ("ouch-ouch disease") was caused by cadmium poisoning.
Cadmium poisoning is usually diagnosed by its symptoms, particularly if there is reason to believe that the patient has been exposed to cadmium. Because patients may not request treatment for up to a day following cadmium exposure, diagnosticians should carefully question any patient who shows symptoms consistent with this condition.
Signs and Symptoms:
Acute exposure to cadmium fumes may cause flu like symptoms including chills, fever, and muscle ache. Symptoms may resolve after a week if there is no respiratory damage. More severe exposures can cause tracheo-bronchitis, pneumonitis, and pulmonary edema. Symptoms of inflammation may start hours after the exposure and include cough, dryness and irritation of the nose and throat, headache, dizziness, weakness, fever, chills, and chest pain.
Inhaling cadmium-laden dust quickly leads to respiratory tract and kidney problems which can be fatal (often from renal failure). Ingestion of any significant amount of cadmium causes immediate poisoning and damage to the liver and the kidneys. Compounds containing cadmium are also carcinogenic.
The bones become soft (osteomalacia), lose bone mass and become weaker (osteoporosis). This causes the pain in the joints and the back, and also increases the risk of fractures. In extreme cases of cadmium poisoning, the mere body weight causes a fracture.
The kidneys lose their function to remove acids from the blood in proximal renal tubular dysfunction. The kidney damage inflicted by cadmium poisoning is irreversible and does not heal over time. The proximal renal tubular dysfunction creates low phosphate levels in the blood (hypophosphatemia), causing muscle weakness and sometimes coma. The dysfunction also causes gout, a form of arthritis due to the accumulation of uric acid crystals in the joints because of high acidity of the blood (hyperuricemia). Another side effect is increased levels of chloride in the blood (hyperchloremia). The kidneys can also shrink up to 30%.
Other patients lose their sense of smell (anosmia).
In the 1950s and 1960s industrial exposure to cadmium was high. But as the toxic effects of cadmium became apparent, industrial limits on cadmium exposure have been reduced in most industrialized nations and many policy makers agree on the need to reduce exposure further. While working with cadmium it is important to do so under a fume hood to protect against dangerous fumes. Silver solder, for example, which contains cadmium, should be handled with care. Serious toxicity problems have resulted from long-term exposure to cadmium plating baths.
Buildup of cadmium levels in the water, air, and soil has been occurring particularly in industrial areas. Environmental exposure to cadmium has been particularly problematic in Japan where many people have consumed rice that was grown in cadmium contaminated irrigation water.
Food is another source of cadmium. Plants may only contain small or moderate amounts in non-industrial areas, but high levels may be found in the liver and kidneys of adult animals.
Cigarettes are also a significant source of cadmium exposure. Although there is generally less cadmium in tobacco than in food, the lungs absorb cadmium more efficiently than the gut. This goes for marijuana as well as tobacco.
Aside from tobacco smokers, people who live near hazardous waste sites or factories that release cadmium into the air have the potential for exposure to cadmium in air. However, numerous state and federal regulations in the United States control the amount of cadmium that can be released to the air from waste sites and incinerators so that properly regulated sites are not hazardous. The general population and people living near hazardous waste sites may be exposed to cadmium in contaminated food, dust, or water from unregulated releases or accidental releases. Numerous regulations and use of pollution controls are enforced to prevent such releases.
Workers can be exposed to cadmium in air from the smelting and refining of metals, or from the air in plants that make cadmium products such as batteries, coatings, or plastics. Workers can also be exposed when soldering or welding metal that contains cadmium. Approximately 512,000 workers in the United States are in environments each year where a cadmium exposure may occur. Regulations that set permissible levels of exposure, however, are enforced to protect workers and to make sure that levels of cadmium in the air are considerably below levels thought to result in harmful effects.
Some sources of phosphate in fertilizers contain Cadmium in amounts of up to 100 mg/kg, which can lead to an increase in the concentration of Cadmium in soil. (for example in New Zealand) Nickel-cadmium batteries are one of the most popular and most common cadmium-based products.
The pathophysiology of the heavy metal toxidromes remains relatively constant. For the most part, heavy metals bind to oxygen, nitrogen, and sulfhydryl groups in proteins, resulting in alterations of enzymatic activity. This affinity of metal species for sulfhydryl groups serves a protective role in heavy metal homeostasis as well. Increased synthesis of metal binding proteins in response to elevated levels of a number of metals is the body´s primary defense against poisoning. For example, the metalloproteins are induced by many metals. These molecules are rich in thiol ligands, which allow high-affinity binding with cadmium, copper, silver, and zinc among other elements. Other proteins involved in both heavy metal transport and excretion through the formation of ligands are ferritin, transferrin, albumin, and hemoglobin.
Although ligand formation is the basis for much of the transport of heavy metals throughout the body, some metals may compete with ionized species such as calcium and zinc to move through membrane channels in the free ionic form.
The study of health effects of cadmium with respect to the cardiovascular system and calcium metabolism disproved the hypothesis that exposure to cadmium would lead to an increase in blood pressure and in the prevalence of hypertension and other cardiovascular diseases. On the other hand, there was a positive relationship between urinary cadmium (Cd-U) and both serum alkaline phosphatase activity and urinary excretion of calcium. The regression coefficients obtained after adjustment for significant co-variates indicated that, when Cd-U increased two-fold, serum alkaline phosphatase and urinary calcium rose by 4% and 0.25 mmol/24 h, respectively. These findings suggest that calcium metabolism is gradually affected as cadmium accumulates in the body. The morbidity associated with the latter phenomenon is still unknown, and requires further investigation, preferably in a longitudinal prospective population study, in which the incidence of morbid events would be monitored in relation to the cadmium body burden.
Cadmium derives its toxicological properties from its chemical similarity to zinc an essential micronutrient for plants, animals and humans. Cadmium is biopersistent and, once absorbed by an organism, remains resident for many years (over decades for humans) although it is eventually excreted.
In humans, long-term exposure is associated with renal disfunction. High exposure can lead to obstructive lung disease and has been linked to lung cancer, although data concerning the latter are difficult to interpret due to compounding factors. Cadmium may also produce bone defects (osteomalacia, osteoporosis) in humans and animals. In addition, the metal can be linked to increased blood pressure and effects on the myocardium in animals, although most human data do not support these findings.
The average daily intake for humans is estimated as 0.15µg from air and 1µg from water. Smoking a packet of 20 cigarettes can lead to the inhalation of around 2-4µg of cadmium, but levels may vary widely.
The organ systems affected and the severity of the toxicity vary with the particular heavy metal involved, the age of the individual, and the level of toxicity.
Also see Signs and Symptoms.
Approximately 2-7% of ingested cadmium is absorbed through the gastrointestinal tract, and its absorption is enhanced when the diet is deficient in calcium, iron, or protein (Lauwerys 1994). Absorption through the respiratory tract is more efficient, ranging from 15% to as much as 50% of an inhaled dose (Lauwerys 1994, Goyer 1996). Both these routes are potential sources of exposure in children.
Cadmium binds to red blood cells, plasma albumin, and metallothionein, which is synthesized in the liver and also by the placenta. Methallothionein may serve as a barrier to protect the fetus (Goyer 1996); however, in cases of excessive maternal exposure, it appears that some cadmium will cross the placenta (Frery 1993).
Cadmium is initially detoxified in the liver through the formation of a metallothionein-cadmium complex, which is slowly released from that organ. Although initially non-toxic, the cadmium-metallothionein complex can be nephrotoxic as it accumulates in the kidneys (Goyer 1996, Dorian 1995).
The average blood level of cadmium in adults without excessive or occupational exposure is about 1 mg/dL or less, as is the amount excreted in the urine in the adult population. Blood and/or urinary cadmium excretions exceeding 5 mg/dL generally indicate excessive exposure (Barnhart 1984). Human breast milk concentrations of cadmium are usually very low.
Cases of severe acute cadmium poisoning are rare.
Heavy metal toxicities are relatively uncommon. However, failure to recognize and treat heavy metal toxicities can result in significant morbidity and mortality.
It affects all races.
Several points are of concern in heavy metal toxicity with respect to age. Generally, children are more susceptible to the toxic effects of the heavy metals and are more prone to accidental exposures. However, company workers, comprising of adults, are no exception.
* Little or no difference in prevalence exists.
* Occupations with heavy metal exposure that predominantly involve a particular sex are associated with higher rates of exposure in that sex.
* A history of ingestion or exposure is the most critical aspect of diagnosing heavy metal toxicity. A complete history, including occupational, hobby, recreational, and environmental exposure is crucial in diagnosing heavy metal toxicity.
* Most acute presentations involve industrial exposure.
* A history of ingestion often leads to the diagnosis in children.
- Metallic taste and increased salivation
- Irritation of respiratory and gastrointestinal tracts.
- Inhalation can cause respiratory toxicity after a latency period of several hours, including a mild, self-limited illness of fever, cough, malaise, headaches and abdominal pain, similar to metal fume fever.
- At higher doses, chemical pneumonitis may occur, with labored breathing, chest pain, and a sometimes fatal hemorrhagic pulmonary edema.
- Ingested cadmium causes nausea, vomiting, diarrhea, abdominal pain and tenesmus.
Chronic symptoms may include:
- Kidney damage (proteinuria and azotemia) anemia, liver injury (jaundice) and defective bond structure. Chronic obstructive pulmonary disease for those chronically exposed by inhalation.
Bioaccumulation of Cadmium in the system.
Three main tests are used for measuring cadmium exposures: cadmium in whole blood, cadmium in urine, and measurement of plasma proteins in urine. A number of other tests may be employed in investigating cadmium-related health effects.
A. Cadmium in Blood
Mainly because of ease of analysis, cadmium in whole blood has been used as a biological indicator of occupational exposures. Cadmium concentrations in blood are mainly a reflection of recent exposure. The ACGIH suggests that monitoring in blood is preferred during the initial year of exposure and whenever changes in the degree of exposure are suspected. In workers not currently exposed, cadmium in blood decreases substantially. When declining blood cadmium levels reach a steady state, they are considered to reflect body burden from previous exposures.
Normal values of cadmium in blood of non-smokers are generally less than 1 ug/l. Higher average values of 1.4 to 4.2 ug/l are found in smokers, though individual blood cadmium levels in smokers may exceed these values.
In 1991, the Ontario Ministry of Labour suggested medical assessment for exposed workers whose blood cadmium level reaches 11 ug/l. However, OSHA recently chose 5 ug Cd/l of whole blood as a level at which further medical surveillance is required of American workers. If this level of cadmium in blood is accompanied by protein in the urine, then OSHA requires workers to be medically removed. A level of 15 ug Cd/l is cause for removal without proteinuria. The ACGIH has also recently proposed a Biological Exposure Index of 5 ug/l of cadmium in blood. The BEI "is intended to prevent the potential for increased urinary excretion of markers of renal dysfunction in almost all workers".
B. Cadmium in Urine
Cadmium concentration in urine is considered to be more reflective of body burden in currently-exposed workers than cadmium in blood, and is the most widely used biological measure of chronic exposure to cadmium. Cadmium in urine increases with age, cigarette smoking, and exposures in the general and occupational environments.
The normal concentration of cadmium in urine is from 0.1 to 1 ug/g creatinine. Until recently, a measure of 10 ug Cd/g creatinine has been regarded as a threshold for kidney effects. However, a number of recent studies have cast doubt on this figure. Evidence of subtle kidney effects are demonstrated at levels a low as 2 ug Cd/g creatinine. Levels of 5 to 1ug Cd/g creatinine are associated with a 10% risk of increased excretion of enzymes and proteins. OSHA recently chose a level of 3 ug Cd/g creatinine as a trigger for enhancing medical surveillance of cadmium. If this level of cadmium in urine is accompanied by proteinuria, then OSHA requires medical removal of the affected worker. A level of 15 ug/g creatinine is cause for removal without proteinuria. The ACGIH has recommended a new Biological Exposure Index (BEI) of 5 ug/g creatinine for cadmium in urine.
C. Markers of early Renal Effects from Cadmium Exposure
While not a measure of cadmium exposure per se, the increase of proteins in urine is a marker of damage to the kidneys which precedes or accompanies most health effects associated with cadmium exposure. One particular protein -- beta2-microglobulin (BMG) -- has been extensively used as an indicator of cadmium-related damage to the proximal tubules of the kidney. BMG excretion may be elevated due to other causes: anti-cancer drugs, antibacterial antibiotics such as streptomycin, anti-inflammatory compounds, myeloma, flu and upper respiratory tract infections. These factors can be readily identified, and need not confound the diagnosis of cadmium-related proteinuria.
Levels of BMG are considered elevated by most investigators at 300 ug/g creatinine, although levels as low as 200 or as high as 500 ug/g creatinine have been suggested as abnormal. OSHA mandates a removal level of 1500 ug BMG/g creatinine, if cadmium levels in blood or urine are elevated.
Exposures to cadmium for 20 years at a level 50 ug/m3 (0.05 mg/m3), the current Ontario limit, give rise to a greatly increased incidence of tubular proteinuria as indicated by output of BMG.
In recent years, a number of other markers for cadmium effects have been recommended. Several of these markers appear to be more sensitive to the early effects of cadmium on the kidney, and/or more stable than BMG in urine. The following markers have been assessed and shown to have significant association with cadmium exposure:
- retinol-binding protein (RBP) in urine;
- albumin in urine;
- N-acetyl-D-glucosaminidase (NAG) in urine;
- metallothionein (MT) in urine;
- urinary transferrin;
- most tubular antigens.
Conventional indicators of renal function such as total urinary protein, serum urea, and serum creatinine are considered insensitive indicators of early renal dysfunction, but may indicate the progression of cadmium-related damage.
Direct Measurement of Cadmium Concentration in Liver and Kidney
Neutron activation analysis is a new method which allows for the direct measurement of the cadmium burden in the liver and kidney. The technique involves use of an ultrasonic scan to precisely locate the target organs, followed by irradiation with a neutron beam which allows assessment of organ burden by measurement of cadmium-specific gamma rays. The radiation dosage is less than most conventional x-rays.
For workers who have been out of exposure for some time or who have suffered kidney damage, this technique can provide a more accurate measure of body burden and may help in determining if non-specific diseases such as emphysema are cadmium-related. This equipment has been employed in England to resolve compensation disputes. Only one Canadian facility is currently equipped to carry out this kind of analysis -- at McMaster University in Hamilton.
Estimation of liver burden is considered more appropriate because once renal damage occurs, cadmium excretion increases and the kidneys lose their cadmium burden. Liver and kidney burdens increase until a 40 ppm concentration is reached in the liver, after which kidney levels decrease while liver burden continues to rise. One study measured a mean liver cadmium burden of 0.6 ppm in non-exposed controls.
X-ray fluorescence, another technique for in vivo measurement of cadmium body burden, has also been developed recently, but is not generally available at this time.
Cadmium Exposure Limits:
Exposure limits for cadmium and its compounds have declined steadily over time, as the knowledge about cadmium-related diseases has grown. The current Ontario time-weighted average limit of 0.05 mg/m3 is considered too high by many scientists. Ontario is currently considering a new limit of 0.02 mg/m3, to match the lower Dutch limit. The Dutch Working Group of Experts which reviewed cadmium's toxicity in 1980, actually proposed a health-based limit of 0.01 mg/m3. However, this limit was not considered feasible at the time. The cadmium limit is again under review in the Netherlands.
In Sweden, a limit of 0.01 mg/m3 is in effect for all new industries employing cadmium or its compounds.
In 1990, the American Conference of Governmental Industrial Hygienists (ACGIH), an influential body which produces a list of Threshold Limit Values adopted by many governments as enforceable occupational exposure limits, proposed a new total dust limit of 0.01 mg/m3 for cadmium, and a 0.002 mg/m3 respirable dust limit.
In 1992, after an exhaustive review of occupational exposures, toxicity and feasibility issues, the U.S. Occupational Safety and Health Administration adopted a new cadmium limit of 0.005 mg/m3, one-tenth of the current Ontario limit, in order to prevent kidney effects and cancer in exposed workers.
In 1991 and 1992, a lively debate occurred in the scientific literature on the question of a protective occupational exposure limit. Proposed protective limits varied from a low of 0.001 mg/m3 through 0.01 mg/m3. One author suggested a limit in the range of 0.01 - 0.1 mg/m3 might be protective, but put some emphasis on 0.02 mg/m3. Although there is no consensus on the level of a protective limit, there is considerable agreement in the scientific community that a limit of 0.05 mg/m3 does not provide sufficient protection, and that kidney effects occur at this level of exposure.
Independent Interventions/First Aid:
NOTE! PREVENT DISPERSION OF DUST! STRICT HYGIENE!
|General First Aid: IN ALL CASES CONSULT A DOCTOR!|
Route of Exposure
|Inhalation||Cough. Headache. Symptoms may be delayed (see Notes).||Fresh air rest. Half-upright position. Artificial respiration if indicated. Refer for medical attention.|
|Skin||Remove contaminated clothes. Rinse and then wash skin with water and soap.|
|Eyes||Redness. Pain.||First rinse with plenty of water for several minutes (remove contact lenses if easily possible) then take to a doctor.|
|Ingestion||Abdominal pain. Diarrhoea. Headache. Nausea. Vomiting.||Rest. Refer for medical attention.|
|Notes for ICSC Information|
|Reacts violently with fire extinguishing agents such as water foam carbon dioxide and halons. Depending on the degree of exposure periodic medical examination is indicated. The symptoms of lung edema often do not become manifest until a few hours have passed and they are aggravated by physical effort. Rest and medical observation are therefore essential. Do NOT take working clothes home.|
If intentional ingestion or overdose is suspected, place patient in closely a monitored unit and consult a medical toxicologist and psychiatrist.
* Contact a certified poison control center or medical toxicologist.
* Consult a gastroenterologist if the possibility of corrosive GI effects is present.
Eat a balanced diet that provides enough calcium, iron, protein, and zinc.
No treatment has been proven effective for cadmium poisoning.
|Follow Up Treatment/Management|
Further inpatient care:
Encourage patient to eat foods rich in calcium, iron, protein, and zinc.
Further outpatient care:
* Care must be taken to remove the source of heavy metal contamination.
* Report industrial-related toxicities to OSHA or; report childhood cases to the local health department.
· Avoid Smoking.
· Avoid hazardous areas or hazardous work such as welding.
Complications include pneumonitis and pulmonary edema. Chronic exposure may cause anemia, emphysema or renal failure, and cadmium may be a risk factor in the development of prostate or lung cancer.
The prognosis depends on the nature and severity of the cadmium load. Most cases of mild exposure resolve spontaneously after a few days. In other cases, cadmium can lead to permanent damage with shortened lifespan, or even death.
Cadmium may be carcinogenic.
Long-term exposure may also result in bone defects including osteoporosis.
Teach the patient on how to prevent cadmium poisoning.
· Failure to report such toxicity to the local health authorities.
Return to List of Diseases and Disorders |
The term interfaith dialogue refers to cooperative, constructive and positive interaction between people of different religious traditions (i.e., "faiths") and/or spiritual or humanistic beliefs, at both the individual and institutional levels. It is distinct from syncretism or alternative religion, in that dialogue often involves promoting understanding between different religions to increase acceptance of others, rather than to synthesize new beliefs.
Throughout the world there are local, regional, national and international interfaith initiatives; many are formally or informally linked and constitute larger networks or federations. The often quoted "There will be no peace among the nations without peace among the religions. There will be no peace among the religions without dialogue among the religions" was formulated by Dr Hans Küng, a Professor of Ecumenical Theology and President of the Foundation for a Global Ethic.
The United States Institute of Peace published works on interfaith dialogue and peacebuilding including a Special Report on Evaluating Interfaith Dialogue Interfaith dialog forms a major role in the study of religion and peacebuilding.
To some, the term interreligious dialogue has the same meaning as interfaith dialogue. Neither are the same as Nondenominational Christianity. The World Council of Churches, though. distinguishes between 'interfaith' and 'interrreligious.' To the WCC, 'interreligious' refers to action between different Christian denominations. So, 'interfaith' refers to interaction between different faith groups such as Muslim and Christian or Hindu and Jew for example.
The history of interfaith dialogue is as ancient as the religions since men and women when not at war with their neighbours have always made an effort to understand them (not least because understanding is a strategy for defence, but also because for as long as there is dialogue wars are delayed). History records many examples of interfaith initiatives and dialogue throughout the ages.
- Interfaith dialogue and action have taken place for many centuries. The Emperor Akbar the Great, for example, encouraged tolerance in Mughal India, a diverse nation with people of various faith backgrounds, including Islam, Hinduism, Sikhism, and Christianity. Religious pluralism can also be observed in other historical contexts, including Muslim Spain. Zarmanochegas (Zarmarus) (Ζαρμανοχηγὰς) was a monk of the Sramana tradition (possibly, but not necessarily a Buddhist) from India who journeyed to Antioch and Athens while Augustus (died 14 CE) was ruling the Roman Emprire.
- The Ottoman Turks' administration of the Balkans from the 15th to 19th centuries provides another historical example of generally peaceful coexistence between peoples of different faiths, including Sufi and non-Sufi Muslims, Roman Catholics, Orthodox Christians and Jews. The tolerant context of this period contrasts dramatically with the ethnic strife and atrocities in the region during the Yugoslav wars of the 1990s.
- There have been several meetings referred to as a Parliament of the World’s Religions, most notably the World's Parliament of Religions of 1893, the first attempt to create a global dialogue of faiths. The event was celebrated by another conference on its centenary in 1993. This led to a new series of conferences under the official title "Parliament of the World's Religions".
- Early 20th Century - dialogue started to take place between the Abrahamic faiths - Christianity, Judaism, Islam and Bahá'í.
- The 1960s - The interfaith movement gathered interest.
- 1965 - The Roman Catholic Church issued the Vatican II document Nostra Aetate, instituting major policy changes in the Catholic Church's policy towards non-Christian religions.
- In the late 1960s interfaith groups such as the Clergy And Laity Concerned (CALC) joined around Civil Rights issues for African-Americans and later were often vocal in their opposition to the Vietnam War.
- September 11, 2001 - After September 11, under the leadership of James Parks Morton, Dean Emeritus of the Cathedral of St. John the Divine, the Interfaith Center of New York's mission became increasingly centered on providing assistance to immigrant and disenfranchised communities whose religious leaders were often the only source of knowledge for new immigrants about coping with a new life in an urban environment like New York City. New programs were launched that responded to the needs of these constituents, combining practical information about establishing civic connections and information about other religions with insight about common social concerns. New programs included Religious Communities and the Courts System (2003), Teacher Education in American Religious Diversity (2003), Mediation for Religious Leaders (2005), and Religious Diversity Training for Social Workers (2005).
- On October 13, 2007 Muslims expanded their message. In A Common Word Between Us and You, 138 Muslim scholars, clerics and intellectuals unanimously came together for the first time since the days of the Prophet[s] to declare the common ground between Christianity and Islam.
- In 2008, through the collaboration of The Hebrew Union College, Omar Foundation, and the University of Southern California Center for Muslim-Jewish Engagement was created. This inter-faith think tank began to hold religious text-study programs throughout Los Angeles and has an extensive amount of resources on its website including scholarly articles about Creationism, Abraham and Human Rights.
- July 2008 - A historic interfaith dialogue conference was initiated by King Abdullah of Saudi Arabia to solve world problems through concord instead of conflict. The conference was attended by religious leaders of different faiths such as Christianity, Judaism, Buddhism, Hinduism, and Taoism and was hosted by King Juan Carlos of Spain in Madrid.
- January 2009, at Gujarat’s Mahuva, the Dalai Lama inaugurated an interfaith "World Religions-Dialogue and Symphony" conference convened by Hindu preacher Morari Bapu from January 6 to 11th 2009. This conference explored ways and means to deal with the discord among major religions, according to Morari Bapu. Participants included Prof. Samdhong Rinpoche on Buddhism, Diwan Saiyad Zainul Abedin Ali Sahib (Ajmer Sharif) on Islam, Dr. Prabalkant Dutt on non-Catholic Christianity, Swami Jayendra Saraswathi on Hinduism and Dastur Dr. Peshtan Hormazadiar Mirza on Zoroastrianism.
- July 2009, the Vancouver School of Theology opened the Iona Pacific: Inter-Religious Centre for Social Action, Research, and Contemplative Practice under the leadership of Principal and Dean, Dr. Wendy Fletcher, and Director, Rabbi Dr. Robert Daum.
Policies of religions to interfaith dialogue
Bahá'í Faith
Interfaith and multi-faith interactivity is integral to the teachings of the Bahá'í Faith. Its founder Bahá'u'lláh enjoined his followers to "consort with the followers of all religions in a spirit of friendliness and fellowship." Bahá'ís are often at the forefront of local inter-faith activities and efforts. Through the Bahá'í International Community agency, the Bahá'ís also participate at a global level in inter-religious dialogue both through and outside of the United Nations processes.
In 2002 the Universal House of Justice, the global governing body of the Bahá'ís, issued a letter to the religious leadership of all faiths in which it identified religious prejudice as one of the last remaining "isms" to be overcome, enjoining such leaders to unite in an effort to root out extreme and divisive religious intolerance.
Buddhism is a religion which teaches people to 'live and let live'. In the history of the world, there is no evidence to show that Buddhists have interfered or done any damage to any other religion in any part of the world for the purpose of introducing their religion. Buddhists do not regard the existence of other religions as a hindrance to worldly progress and peace.
The 14th century Zen master Gasan Joseki indicated that the Gospels were written by an enlightened being:
- "And why take ye thought for raiment? Consider the lilies of the field, how they grow. They toil not, neither do they spin, and yet I say unto you that even Solomon in all his glory was not arrayed like one of these...Take therefore no thought for the morrow, for the morrow shall take thought for the things of itself."
- Gasan said: "Whoever uttered those words I consider an enlightened man."
The 14th Dalai Lama has done a great deal of interfaith work throughout his life. He believes that the "common aim of all religions, an aim that everyone must try to find, is to foster tolerance, altruism and love". He met with Pope Paul VI at the Vatican in 1973. He met with Pope John Paul II in 1980 and also later in 1982, 1986, 1988, 1990, and 2003. During 1990, he met in Dharamsala with a delegation of Jewish teachers for an extensive interfaith dialogue. He has since visited Israel three times and met during 2006 with the Chief Rabbi of Israel. In 2006, he met privately with Pope Benedict XVI. He has also met the late Archbishop of Canterbury Dr. Robert Runcie, and other leaders of the Anglican Church in London, Gordon B. Hinckley, late President of the Church of Jesus Christ of Latter-day Saints (Mormons), as well as senior Eastern Orthodox Church, Muslim, Hindu, Jewish, and Sikh officials.
In 2010, the Dalai Lama was joined by Rev. Katharine Jefferts Schori, presiding bishop of the Episcopal Church, Chief Rabbi Lord Jonathan Sacks of the United Hebrew Congregations of the Commonwealth, and Islamic scholar Professor Seyyed Hossein Nasr of George Washington University when Emory University's Center for the Study of Law and Religion hosted a "Summit on Happiness".
|This section does not cite any references or sources. (April 2013)|
Traditional Christian doctrine is Christocentric, meaning that Christ is held to be the sole full and true revelation of the will of God for humanity. In a Christocentric view, the elements of truth in other religions are understood in relation to the fullness of truth found in Christ. God is nevertheless understood to be free of human constructions. Therefore, God the Holy Spirit is understood as the power who guides non-Christians in their search for truth, which is held to be a search for the mind of Christ, even if "anonymously," in the phrase of Catholic theologian Karl Rahner. For those who support this view, anonymous Christians belong to Christ now and forever and lead a life fit for Jesus' commandment to love, even though they never explicitly understand the meaning of their life in Christian terms.
While the conciliar document Nostra Aetate has fostered widespread dialogue, the declaration Dominus Iesus nevertheless reaffirms the centrality of the person of Jesus Christ in the spiritual and cultural identity of Christians, rejecting various forms of syncretism.
Pope John Paul II was a major advocate of interfaith dialogue, promoting meetings in Assisi in the 1980s. Pope Benedict XVI has taken a more moderate and cautious approach, stressing the need for intercultural dialogue, but reasserting Christian theological identity in the revelation of Jesus of Nazareth in a book published with Marcello Pera in 2004.
For traditional Christian doctrine, the value of inter-religious dialogue is confined to acts of love and understanding toward others either as anonymous Christians or as potential converts.
In mainline liberal Protestant traditions, however, as well as in the emerging church, these doctrinal constraints have largely been cast off. Many theologians, pastors, and lay people from these traditions do not hold to uniquely Christocentric understandings of how God was in Christ. They engage deeply in interfaith dialogue as learners, not converters, and desire to celebrate as fully as possible the many paths to God.
Much focus in Christian interfaith dialogue has been put on Christian–Jewish reconciliation. One of the oldest successful dialogues between Jews and Christians has been taking place in Mobile, Alabama. It began in the wake of the call of the Second Vatican Council (1962–1965) of the Roman Catholic Church for increased understanding between Christians and Jews. The organization has recently moved its center of activity to Spring Hill College, a Catholic, Jesuit institution of higher learning located in Mobile. Reconciliation has been successful on many levels, but has been somewhat complicated by the Arab-Israeli conflict in the Middle East, where a significant minority of Arabs are Christian.
Orthodox Judaism forbids interfaith dialogue. Prominent Rabbinic authorities in ruling on this issue, further caution that the intent is to convert Jews, that participation leads to confusion and wrong ideas and that Judaism's prohibition of proselytism, combined with other religions' "missionary zeal", creates an unbalanced dynamic such that the "dialogue" effectively becomes a monologue. The Modern Orthodox movement allows narrow exchanges on social issues while forbidding discussion of doctrine. Reform Judaism, Reconstructionist Judaism and Conservative Judaism encourage interfaith dialogue.
||This article needs additional citations for verification. (September 2010)|
Islam has long encouraged dialogue to reach truth (and not interfaith dialogue which seeks to find common between people and leave differences aside). Islam also stressed that the supreme law of the land should be Islam and that Islam regulates all life affairs and therefore regulates how non Muslim and Muslims live under an Islamic state, with historical examples coming from Muslim Spain, Mughal India, and even starting as far back as Muhammad's time, where people of the Abrahamic Faiths lived in harmony.
Many traditional and religious texts and customs of the faith have encouraged this, including specific verses in the Quran, such as: "O people! Behold, we have created you from a male and a female and have made you into nations and tribes so that you might come to know one another. Verily, the noblest of you in the sight of God is the one who is most deeply conscious of Him. Behold, God is all-knowing, all-aware." [Qur'an 49:13]
In recent times, Muslim theologians have advocated inter-faith dialogue on a large scale, something which is new in a political sense. The declaration A Common Word of 2007 was a public first in Christian-Islam relations, trying to work out a moral common ground on many social issues.
Relations between Muslims and Jews remain quite difficult, exacerbated by the Israeli-Palestinian conflict. There are inter-Muslim issues in between Sunnis and Shiites that are very much unresolved in the Middle East. Also, relations between Muslims and Hindus in India and Pakistan could theoretically be much better if interfaith efforts were more successful.
According to the Ahmadiyya understanding, interfaith dialogues are an integral part of developing inter-religious peace and the establishment of peace. The Ahmadiyya Community has been organising interfaith events locally and nationally in various parts of the world in order to develop a better atmosphere of love and understanding between faiths. Various speakers are invited to deliver a talk on how peace can be established from their own or religious perspectives.
Zoroastrianism has long encouraged interfaith, all the way from Cyrus the Great's speech in Babylon, which permitted the population to keep following their own religion and keep speaking their own language. Cyrus did not enforce the state religion unto the people. As well, Cyrus freed all the Jewish slaves from Babylon, which earned him a place in the Jewish scriptures. Zoroastrians believe that all religions are equal, and that their religion is not superior to other religions. They believed that the Prophet Zoroaster implied the religion unto them, and did not convert each of them. Therefore, they do not even accept converts into their religion. All adherents must be born into the religion.
Interfaith organisations
Muslim Christian Dialog Forum is an Interfaith Dialog Forum established by Minhaj-ul-Quran-a sufi based moderate Islamic Organization started by Tahir_ul_Qadri to promote religious tolerance and cultural co-existence between Muslims and Christian community.
Interfaith Encounter Association (IEA) was established in 2001 and works to build genuine coexistence and sustainable peace, through joint community building on the grassroots level, using interactive interfaith dialogue as its vehicle. The a-political and all-inclusive approach of the organization and its activities continuously form the human infrastructure for peace in the Holy Land and the Middle East. In its ten years of existence, the IEA have held – in its three regional focuses: in Israel, between Israelis and Palestinians and in the larger Middle East – more than 1000 programs, with thousands of participants. A most significant fact is that the participants in IEA programs include people of all political and religious views, as well as all ages, genders, walks of life etc.; and that the vast majority of them have met 'the other' for the first time through IEA. The IEA have formed till now 41 on-going community-groups of interfaith encounter – from the Upper Galilee to Eilat, including 10 groups that bring together on a regular basis Israelis and West Bank Palestinians. Among the latter we maintain the three only groups in the country that bring together Palestinians with Settlers. IEA maintains working relations with 7 Palestinian organizations, across the West Bank and the Gaza Strip and is a founding partner of the Middle East Abrahamic Forum, with additional organizations from Egypt, Iran, Jordan, Lebanon, Morocco, Tunisia and Turkey.
Messiah Foundation International is an interfaith organisation which aims to promote mutual love, peace and understanding between members of all religions and faiths through the spiritual sciences taught by Ra Gohar Shahi. MFI has centres across the globe, including in the United States of America, Canada, the United Kingdom of Great Britain, India, Pakistan, Sri Lanka, and Japan.
Project Interfaith is a non-profit organization that aims to grow understanding, respect and relationships among people of all faiths, beliefs and cultures. The goals of the organization are approached through online media resources (particularly RavelUnravel) as well as community-building programs that educate and engage a variety of audiences on issues of faith, religion, identity, interfaith relations, and religious and cultural diversity.
While there are many essentially religious organisations geared towards working on interfaith issues (see Interreligious organisations) there is also a less common attempt by some governmental institutions to specifically address the diversity of religions (see Australasian Police Multicultural Advisory Bureau for one example).
In India, many organizations have been involved in interfaith activities because of the diversity of religion in the nation.
United Religions Initiative (URI) was founded in 2000 to promote daily, lasting interfaith cooperation, end religiously motivated violence, and create cultures of peace, justice, and healing for the Earth and all living beings. With hundreds of thousands of members in 80+ countries representing over 200 religions and indigenous traditions, URI uses "cooperation circles" to promote dialogue and action.
The Jordanian Interfaith Coexistence Research Center is a Jordanian non-governmental organization for promoting peaceful religious coexistence. It fosters grassroots interfaith dialogue and works on creating interreligious harmony.
The Global Peace Pioneers - GPP is a Pakistan based non-governmental organization advocating for peaceful religious coexistence. It cultivates the seed for grassroots interfaith dialogue and works on creating interreligious harmony with in Pakistan.
United Nations support
On December 2, 2008, Anwarul Karim Chowdhury said:
- "Interfaith dialogue is absolutely essential, relevant, and necessary. ... If 2009 is to truly be the Year of Interfaith Cooperation, the U.N. urgently needs to appoint an interfaith representative at a senior level in the Secretariat."
- The Republic of the Philippines will host a Special Non-Aligned Movement Ministerial Meeting on Interfaith Dialogue and Cooperation for Peace and Development from March 16 to 18 in Manila. During the meeting, to be attended by ministers of foreign affairs of the NAM member countries, a declaration in support of interfaith dialogue initiatives will be adopted. An accompanying event will involve civil society activities.
- In 2010, HM King Abdullah II addressed the 65th UN General Assembly and proposed the idea for a ‘World Interfaith Harmony Week’ to further broaden his goals of faith-driven world harmony by extending his call beyond the Muslim and Christian community to include people of all beliefs, those with no set religious beliefs as well. A few weeks later, HRH Prince Ghazi bin Muhammad presented the proposal to the UN General Assembly, where it was adopted unanimously as a UN Observance Event.
The first week of February, every year, has been declared a UN World Interfaith Harmony Week. The Royal Islamic Strategic Studies Centre released a document which summarises the key events leading up to the UN resolution as well as documenting some Letters of Support and Events held in honour of the week.
Criticism of interfaith dialogue
Conversely, organisations labelled as extremist have been accused of adopting interfaith dialogue as a political front, as well as to raise funds. One commentator has noted of Islamist groups, that, "Interfaith is the perfect ‘do-good’ agenda with which to legitimise their reputation and obfuscate their genuine, more sinister, intentions."
British MP Paul Goodman has questioned the UK Government's decision to fund Campusalam, a University interfaith group, that has received under half a million pounds of taxpayers' money, despite the group's open links to the Lokahi Foundation, widely considered to be an Islamist organisation.
See also
- A Common Word Between Us and You
- Buddhism and Jainism
- Buddhism and Hinduism
- Buddhism and Christianity
- Centre for Dialogue
- Daughters of Abraham
- Ecumenism (Christian)
- Esalen Institute
- The Elijah Interfaith Institute
- Fethullah Gülen
- Gülen movement
- Interfaith Center of New York
- International Center for Religion & Diplomacy
- Jewish views of religious pluralism
- Jordanian Interfaith Coexistence Research Center
- Multifaith space
- Parliament of the World's Religions
- Pontifical Council for Interreligious Dialogue
- Prince Alwaleed Center for Muslim–Christian Understanding
- Relations between Catholicism and Judaism
- Religions for Peace
- Roland de Corneille
- Saltley Gate Peace Group
- Scriptural Reasoning
- Seventh-day Adventist interfaith relations
- Temple of Understanding
- Temple of All Religions
- United Religions Initiative
- United States Institute of Peace
- World Council of Churches
- World Interfaith Harmony Week
- Musser, D & Sunderland, D., War or Words: Interreligious Dialogue as an Instrument of Peace Cleveland: The Pilgrim Press, (2005) 1.
- Smock, D. (ed), (2002)Interfaith Dialogue and Peacebuilding Washington, DC: US Institute of Peace Press
- Abu Nimer, M., et al, (2007) Unity in Diversity: Interfaith Dialogue in the Middle East Washington, DC: US Institute of Peace Press
- Renee Garfinkel,What Works: Evaluating Interfaith Dialogue, United State Institute of Peace, Special Report #123, (2004)
- "Interreligious Dialogue". World Council of Chruches. Retrieved 24 July 2012.
- Progressive Scottish Muslims: Learning Interfaith from the Mughals: Akbar the Great (1556-1605)
- Strabo, xv, 1, on the immolation of the Sramana in Athens (Paragraph 73).
- Dio Cassius, liv, 9.
- Give Peace a Chance: Exploring the Vietnam Antiwar Movement ISBN 978-0-8156-2559-9
- http://www.saudiembassy.net/press-releases/press07170801.aspx Saudi Embassy - Saudi King Abdullah Commences Interfaith Dialogue Conference in Madrid, Spain
- Saudi Gazette - Let concord replace conflict – Abdullah
- Dalai Lama inaugurates 6-day world religions meet at Mahua
- Dalai Lama to inaugurate inter-faith conference
- Tablets of Bahá'u'lláh, page 22, Bahá'u'lláh, From the "Bishárát" (Glad-Tidings).
- Catharine Cookson, ed. (2003). Encyclopedia of religious freedom. Taylor & Francis. p. 9.
- The Buddhist View toward Other Religions
- The Buddhist Attitude Towards Other Religions
- 101 Zen Stories; #16
- Tibetan Buddhism
- Kamenetz,Rodger (1994)The Jew in the Lotus Harper Collins: 1994.
- Top 10 Things Religious Leaders Say about Happiness
- Chabad.info sighting Rambam, Rav Feinstein, Lubavitcher Rebbe
- L'Chaim: 672: Beha'aloscha
- L'Chaim: 673: Sh'lach
- A Modern Orthodox Approach to Interfaith Dialogue
- "Ahmadiyya Muslim Community to hold Peace Conference in Malta". Ahmadiyya times. Retrieved 19 October 2010.
- About the Meeting
- How It Began
- The First UN World Interfaith Harmony Week Booklet
- The Inevitability of Clash of Civilisation |
versión impresa ISSN 1870-3453
Rev. Mex. Biodiv. vol.82 no.4 México dic. 2011
Effects of a snowstorm event on the interactions between plants and hummingbirds: fast recovery of spatiotemporal patterns
Efecto de una tormenta de nieve sobre la interacción colibríplanta: los patrones espaciotemporales se recobran rápido
Román DíazValenzuela1,2 and Raúl OrtizPulido1,3*
1 Laboratorio de Ecología de Poblaciones, Centro de Investigaciones Biológicas, Universidad Autónoma del Estado de Hidalgo. Apartado postal 69, 42001 Pachuca, Hidalgo, México.
2 Laboratorio de Vertebrados, Centro Iberoamericano de la Biodiversidad, Universidad de Alicante, 03080 Alicante, España.
3 Current address: Laboratorio de Ecología de la Conducta, Centro Tlaxcala de Biología de la Conducta, Universidad Autónoma de TlaxcalaUNAM, Carretera TlaxcalaPuebla km 1.5, Colonia Xicohtencatl S/N, 90070 Tlaxcala, Tlaxcala, México. *firstname.lastname@example.org.
Recibido: 25 octubre 2010;
aceptado: 02 marzo 2011
The global climatic change could cause, in some places, appearance of meteorological phenomena considered rare. If we understand the effect of these phenomena on birds we can understand how birds respond to weather changes. We report here the effect of a severe snowfall on hummingbird activity, flower abundance and hummingbirdplant interaction in a temperate forest of central Mexico. During our study we registered 1 hummingbird species (Hylocharis leucotis) and 7 plant species (Fuchsia thymifolia, F. microphyla, Salvia amarissima, S. elegants, Cestrum roseum, Penstemon campanulatus and Lonicera mexicana). Before the sudden climatic phenomena we registered 66 records of hummingbirds, 8 700 flowers, and 6 hummingbird visits to flowers. During the phenomena, there were zero hummingbird records, 160 flowers and zero visits. A month after the event there were 67 hummingbirds records, 1 825 flowers and 13 visits. Hummingbird activity recovered rapidly after the snowstorm, but 6 of 7 plants species lost all their flowers, except for L. mexicana, which received all hummingbird visits a month after the climatic event.
Key words: bird pollinated flowers, global climate change, hummingbirds, ornithophylous plant species, snowstorm effect, sudden changes in the average state of the time.
El cambio climático global podría causar fenómenos meteorológicos considerados como raros. Si entendemos el efecto de estos fenómenos en las aves podríamos comprender como las aves responden a este tipo de cambios en el tiempo. Aquí documentamos el efecto de una nevada severa sobre la actividad de colibríes, la abundancia floral y la interación colibríplanta en un bosque templado del centro de México. Durante nuestro estudio registramos 1 especie de colibrí (Hylocharis leucotis) y 7 especies de plantas (Fuchsia thymifolia, F. microphyla, Salvia amarissima, S. elegants, Cestrum roseum, Penstemon campanulatus y Lonicera mexicana). Antes de la nevada registramos 66 avistamientos de colibríes, 8 700 flores y 6 visitas de colibríes a las flores. Durante el fenómeno hubo cero registros de colibríes, 160 flores y cero visitas. Un mes después del evento climático hubo 67 registros de colibríes, 1 825 flores y 13 visitas. La actividad de los colibríes se recuperó rápido después de la nevada, pero 6 de 7 especies de plantas perdieron todas sus flores, exceptuando a L. mexicana, que recibió todas las visitas de colibríes 1 mes después del evento climático inusual.
Palabras clave: flores polinizadas por colibríes, cambio climático global, colibríes, plantas ornitofílicas, efecto de nevada, cambios repentinos en el estado del tiempo.
The Earth's climatic system is changing (Peterson et al., 2002; Watkinson et al., 2004). The evidence suggests that it could be a global climatic change (Ordano, 2003; Watkinson et al., 2004), and among the consequences are sudden changes in the average state of the weather (SuCAST) which in turn will cause apparition of meteorological phenomena considered rare in some places. Among these SuCAST will be droughts, low temperatures, hurricanes and severe snowstorms (e.g., Watkinson and Gill, 2002; Watkinson et al., 2004). The effect of these phenomena on the fauna could be long or short term. A drought, for example, can affect a locality for months or years, while a hurricane or a severe snowstorm can have immediate catastrophic effects (Rotenberry et al., 1993).
There is evidence that birds are sensitive to SuCAST (e.g., Kalela, 1949; Knopf and Sedgwick, 1987; DeSante et al., 1993; Peterson et al., 2001; TejedaCruz and Sutherland, 2005). However, SuCAST has been only hypothetically related to fauna (e.g., Wagner, 1946; Inouye et al., 1991; Gass et al., 1999). SuCAST, depending on their intensity, duration and affected area, can influence in different ways avian populations (e.g., Pickett and White, 1985; Rotenberry et al., 1993).
The hummingbirdplant interaction is an excellent model to understand how SuCAST influence birds, both directly, through mortality due to changes in temperature, and indirectly, through loss of food sources. Both SuCAST effects has been hypothesized, but never tested in field (Wagner, 1946; Inouye et al., 1991; Gass et al., 1999). The hummingbirdplant interaction could be affected by such events since the abundance and distribution of hummingbirds dependent on the flower / nectar availability in the environment (Stiles, 1977, 1978; OrtizPulido and Díaz, 2001; OrtizPulido and VargasLicona, 2008) and SuCAST affects such flower availability (Wagner, 1946; Gass et al., 1999).
Hummingbirds pollinate close to 15% of angiosperm plants in many vegetation types in America (Buzato et al., 2000). Therefore, a SuCAST influencing hummingbirds will affect the reproductive success of the plants visited by them, and in turn will affect other animals that depend on flowers and fruits.
Here we report how a SuCAST affected the activity of a hummingbird species, the abundance of flowers visited by this species and the interaction between these in a temperate forest located in central Mexico. We present data on hummingbird activity, flower abundance, and interaction between both groups, before, during, an after the SuCAST.
Materials and methods
This study was conducted in the El Chico National Park located near Pachuca City, Hidalgo, Mexico (98°42'33''W, 20°11'22''N, 2 9003 080 m a.s.l.). The climate is semicold subhumid with summer rains (COEDEUAEH, 2004), with a mean annual precipitation of 1 030 mm, mean annual temperature of 14.3°C, and mean temperature of the coldest month of 12.1°C (National Water Commission, unpublished data). The place is dominated by a dense temperate forest, with a 2040 m canopy, and 2 lower vegetation strata: a scarce herbaceous, and a shrub cover (Calderón and Rzedowski, 2001). In the El Chico National Park there are reports of 7 hummingbird species and at least 20 plant species whose flowers are visited by them (Villada, 1873; Arregui, 2004; Mauricio, 2005; MartínezGarcía, 2006; DíazValenzuela, 2008; OrtizPulido and VargasLicona, 2008; OrtizPulido et al., 2008).
On January 1, 2008, there was an intense snowstorm, accompanied by wind and rain (MotaLópez, 2008). It left an accumulated snow layer of approximately 20 cm deep. The temperature varied between 6° and 3°C during 3 days. This type of phenomena is uncommon in the place. During the last 50 years it has been reported only 3 times (1978, 1997, and 2008); however for the first 2 events there is no technical information about their magnitude (Sergio Alarcón Martínez, National Water Commission Hidalgo, com. pers.). During the 2008 snowstorm, close to 50% of the trees (Abies religiosa, Pinaceae) lost many of their branches in a surveyed area of 12 ha (RVD, pers. obs.). The snow remained in the area for at least 15 days.
The only hummingbird species recorded during the 2008 SuCAST was the Whiteeared Hummingbird (Hylocharis leucotis), which is a resident species. At the locality, this species has been recorded foraging throughout the year on flowers of Castilleja tenuiflora, C. moranensis, Penstemon barbatus, P. roseus, P. campanulatus (Fam. Scrophulariaceae), Fuchsia thymifolia, F. microphylla (Onagraceae); Macromeria pringlei (Boraginaceae), Salvia elegans, S. amarissima, Stachys coccínea, Stachys sp., Scutellaria caerulea, Prunella vulgaris (Labiateae), Solanum nigrescens, Cestrum roseum (Solanaceae), Lonicera mexicana (Caprifoliacea), Tillansdia violacea, T. erubescens (Bromeliaceae) and Senecio angulifolius (Compositae) (Arregui, 2004; DíazValenzuela, 2008; OrtizPulido et al., 2008).
Our fieldwork was carried out before, during, an after the SuCAST, from December 2007 to February 2008. We established a 350 x 350 m plot (i.e. 12 ha) on the north face of a hill. The area was divided into 49 subplots of 50 x 50 m. We registered all visual and acoustical hummingbird records (from here, hummingbird activity), and introductions of hummingbird bills into flowers (from here, hummingbird visits) in a 25 m radius from the center of every subplot. Every subplot was surveyed from 07:00 to 12:00 h during 10 minutes every month, registering hour, hummingbird species, and hummingbird activity. Each subplot was visited 1 time during the SuCAST. To determine the avian species we used hummingbird field guides (Williamson, 2001; Howell, 2002; OrtizPulido et al., 2006). To determine flower abundance we registered the number of flowers in 4 areas of 9 x 9 m (i.e. 324 m2) in every subplot. The 4 values were averaged to obtain flower abundance per subplot. This technique has proved to indicate a reliable measure for flower abundance in subplots of 50 x 50 m (DíazValenzuela, 2008; OrtizPulido and VargasLicona, 2008). We only took into account open and mature flowers from those plant species whose flowers were visited by hummingbirds during the study period. Plant species identity was determined using a reference collection located in the Herbarium of the Universidad Autonoma del Estado de Hidalgo.
The month long relationship between hummingbird activity and flower abundance, previous, during, and after the SuCAST, was tested using a Spearman correlation test. Monthly relationships were obtained taking into account 49 samples, 1 for each subplot. Spatial autocorrelation in the response variable was tested using the Mantel test provided by the XLstat program (Addinsoft, 2007), and when spatial dependence was found, we used the Spearman correlation test modified by Dutilleul (1993). This test changes the degree of freedom (using decimal places) to achieve the best interpretation of the statistical significance of spatially explicit data (Dutilleul, 1993; Legendre et al., 2002; Rangel et al., 2006).
A month before and after the 2008 SuCAST, we recorded 1 hummingbird species (Whiteeared Hummingbird), but none during the phenomena (the 15 days with snow on the ground). We recorded 66 hummingbird sightings before the SuCAST, zero during the climatic phenomena, and 67 afterward, for a total of 133 hummingbird sightings. From this total we were unable to identify the hummingbird species in 17 instances.
Before, during, and after the SuCAST the richness of flourishing plant species varied (Table 1), as well as the number of openmature flowers. Before SuCAST we recorded 8 700 flowers, 160 during the phenomena, and 1 825 afterward. During SuCAST all the herbaceous flowers disappeared, except for L. mexicana that displayed 160 flowers (Table 1). Apparently, the SuCAST event determined the disappearance of ~98% of the flowers.
We recorded 6 hummingbird visits to flowers (4 to F. microphylla, 1 to S. elegant and 1 to S. amarissima) before SuCAST, zero during the phenomena and 13 visits (all to L. mexicana) afterward. Before and after SuCAST there were significant correlations (Fig. 1) between hummingbird activity and flower abundance (in December rs=0.56, d.f.= 45.70, F= 21.12, p= 0.00003; in February rs= 0.40, df.=51.5, F=9.70, p= 0.00301). During the SuCAST we did not apply a correlation test, because there was no hummingbird activity in the 49 subplots.
The 2008 SuCAST event was related to low levels of hummingbird activity in the El Chico National Park. This effect was most likely caused by hummingbirds dying due to the freezing conditions, and indirectly by a reduction of hummingbird food sources, that forced hummingbirds to enter into torpor, or move to other locations. Before and after the 2008 SuCAST there was a significant correlation between hummingbird activity and flower abundance. This relationship has been reported in other studies without a SuCAST context (Wagner, 1946; Wolf, 1970; Feinsinger, 1976; Stiles, 1977, 1978; Feinsinger and Colwell, 1978; Baltosser, 1989; Inouye et al., 1991; Gutiérrez and Rojas, 2001; OrtizPulido and Díaz, 2001; Lara, 2006; DíazValenzuela, 2008; OrtizPulido and VargasLicona, 2008).
Exactly how the hummingbirds responded to the phenomena we could not reliably determine. They could have suffered an increase in mortality (as suggested by Wagner, 1946; Inouye et al., 1991), made altitudinal movements to search for food (as suggested by Ornelas and Arizmendi, 1995), or employ torpor (as suggested by Bench et al., 1997). We did not record dead hummingbirds in the sampled 12 ha during the SuCAST (but we recorded 4 dead bats in the snow), did not find any evidence of altitudinal movements or detect torpor for the hummingbirds. At nearby sites of El Chico National Park at lower altitudes (approx 2 3502 900 m a. s. l, with 3 different kinds of vegetation: xerophic shrub, Juniperus forest and oak forest) we did not register unusual variation in hummingbird relative abundance. In the lower elevation sites we conducted the same fieldwork effort reported here, and the SuCAST effects were not so severe, with hummingbird activity and flower abundance being similar to other years.
Even though the 2008 SuCAST had an immediate drastic effect on the number of hummingbird sightings, the effect remained shortterm. A month after the SuCAST hummingbird activity returned to the same levels (66 previous vs. 67 after). However, the effect on the plants was more destructive. Only 1 of 7 plant species (L. mexicana) displayed flowers after the climatic phenomenon. Therefore the SuCAST was costly in reproductive terms for the majority of the local plant species. Nevertheless, L. mexicana may have benefited from the SuCAST event. It attained all the hummingbird visits after the SuCAST and perhaps profited from an increase in pollen transference. However this was not tested. All our data indicate that the planthummingbird interaction system underwent a rapid recovery and that L. mexicana was a main element providing nectar sources to hummingbirds a month after the SuCAST.
The rapid recovery detected of the planthummingbird interaction could be explained by the high temporal and spatial variability of the hummingbirdplant interaction systems (DíazValenzuela, 2008; OrtizPulido y VargasLicona, 2008). In nature the flower abundance varies and is an aggregated source in time and space, usually located in isolated patches with high local abundance (e.g. Baltosser, 1989; Fleming, 1992; DíazValenzuela, 2008; OrtizPulido y VargasLicona, 2008). At our study site, without the occurrence of SuCAST, it has been detected that the flower abundance in 50 x 50 m patches can change from thousands to zero within 4 days (OrtizPulido and VargasLicona, 2008). From a local perspective, i.e., 50 x 50 m, it is possible that the change in floral density originated by the 2008 SuCAST is very similar to the variation which occurs naturally in our study site. However, from a landscape perspective, i.e., 12 ha or more, the 2008 SuCAST imposed a serious energetic challenge to hummingbirds, since the flowers of many plant species disappeared from the landscape.
In the 2008 SuCAST event, we detected that flower availability was reduced significantly in the landscape and this was related to hummingbird activity. During normal conditions, the landscape of El Chico National Park shows ~5.68 Kj/ha in floral nectar (OrtizPulido and Lara, unpub. data), and a Whiteeared Hummingbird individual (~3.5 g) needs ~27 kj/day to survive (based on Baltoser, 1989, and Montgomerie y Gass, 1981, estimates). In a SuCAST event, when ~98% of the flowers disappear, the available nectar energy is reduced at the landscape level. In such conditions hummingbirds either die or deal with the stressful conditions, by moving away or using torpor (as suggested by Wagner, 1946; Inouye et al., 1991; Ornelas and Arizmendi, 1995; Bench et al., 1997). The consequences at the community level are different in each case. If hummingbirds die, some community dynamics could be disrupted or will be decreased for some time. Among the disrupted community dynamics are the hummingbirdplant interactions and the seed production of hummingbird pollinated plant species. If the hummingbirds died, we hypothesized that the interactions would be reestablished slowly after the SuCAST, since hummingbird species must explore and colonize the site again. Community dynamics restart would depend, in this case, upon the response of the hummingbirds to the SuCAST. On the other hand, if hummingbirds flew away or used torpor, we hypothesized that the interaction could be restarted swiftly after the SuCAST, as individual hummingbirds would only need to return to the site and continue with their normal activities when the flowers reappear. However, there is evidence that suggests that a bird species can survive only few hours in torpor state (Bicudo et al., 2002), so torpor may only be a temporal answer for hummingbirds to SuCAST events. If hummingbirds fly away or use torpor, community dynamics restart depends on the reaction of the plants to the SuCAST, because hummingbirds are still present in the landscape.
When a SuCAST occurs, surely some hummingbirds die, others fly away and others may endure part of the SuCAST using torpor. Unfortunately, the hummingbird primary reaction to SuCAST has not been determined yet. This study revealed evidence that indicates that hummingbird activity was reduced during a SuCAST event, but the causes remain yet unknown.
We thank the El Chico National Park Direction for permissions and their logistical support in field, Nadja Weisshaupt for the English translation and 2 anonymous reviewers, whose kind suggestions helped us to improve an early version of this paper. Financial support was given by the Programa de Mejoramiento del Profesorado (Secretaria de Educación Pública México), Consejo Nacional de Ciencia y Tecnología (doctoral scholarship to RDV and sabbatical scholarship to ROP) and the Consejo de Ciencia y TecnologíaHidalgo (project Diversidad Biológica del Estado de Hidalgo, segunda fase, FomixCONACyT200895828).
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ATHANASIUS I, the Apostolic Saint, twentieth patriarch of the See of Saint Mark (326-373). Athanasius' life has been treated in detail by numerous authors. These sources can be categorized as follows: (1) the writings of Athanasius himself, which should be considered the most authentic of the sources. These include his historical tracts, encyclicals, an apology to Constantine, another apology against the Arians, his letters to Serapion and to the monks, and his festal letters; (2) the works of contemporary church fathers, including Hilary of Poitiers, BASIL THE GREAT, GREGORY OF NAZIANZUS, and Epiphanius; and (3) chronicles of older historians such as RUFINUS, SOCRATES, Sulpicius Severus, THEODORET, and SOZOMEN, whose authority on details must be taken with some caution. To these may be added the Arabic life
rendered by E. RENAUDOT but originally prepared for the pious Copts, which is simply a legendary account of little historical import. Of course, the official record of the church is represented in the Copto-Arabic SYNAXARION and cannot be overlooked.
The secondary literature on the great saint is profuse, and only a selection of the most prominent biographers may here be mentioned by way of introduction: B. Montfaucon, L. S. de Tillemont, J. A. Moehler, S. Cave, H. G. Opitz, E. Schwartz, L. Atzberger, H. M. Gwatkin, F. L. Cross, and G. Mueller.
Athanasius was probably born in Alexandria around the year 296, although, according to an Arabic document found in DAYR ANBA MAQAR, it is said that his parents originally came from the city of al-Balyana in Upper Egypt. It is possible that his early education took place in the CATECHETICAL SCHOOL OF ALEXANDRIA; it is also possible that he could have attended classes in the Museon where he became conversant with Neoplatonism. As a young man, he must have witnessed the later period of the age of persecutions, though he would have been far too young to recollect incidents related to Maximian's persecution of 303. After Constantine declared Christianity to be the religion of the state, in the Edict of Milan in 312, his family must have suffered through the nascent Arian heresy, a movement destined to be the focal point of his struggle throughout his life. Rufinus and subsequent historians relate a story about Athanasius' boyhood. It is said that Patriarch ALEXANDER I, watching the seashore from his window, saw a group of children playing at Christian baptism; one of the boys played the bishop. Intrigued by this sight, the patriarch summoned the children to his presence and recognized the authenticity of the baptism thus performed. He kept at his court the boy-bishop, Athanasius, who ultimately became his secretary and his closest companion.
The Council of NICAEA in 325 marked the inauguration of the ecumenical movement. The young Athanasius, as Alexander's secretary, was the power behind the throne, and his influence was felt in the composition of the Nicene Creed. Athanasius succeeded Alexander in 326. The new archbishop now faced alone the spreading doctrine of Arianism.
Arius was probably of Libyan origin and a pupil of Lucian of Antioch. He was first ordained by Achillas (d. 311) as presbyter of the important church of Bucalis in Alexandria. An eloquent speaker and a pragmatic thinker, Arius captivated a large congregation in Alexandria with his ideas. He denied the coequality and coeternity of Jesus with the Father and held that the Father created the Son from nothing, only in turn to create the world. Thus the consubstantiality of the Father and the Son was denied by Arius, whose position was supported by Eusebius, bishop of Nicomedia. The idea was rejected by Alexander, and its vehement opponent was Athanasius, who defended his view at the Council of Nicaea using the famous term HOMOOUSION to describe the consubstantiality of the Father and the Son. The defeat of the Arians at that council did not end the controversy nor did it eliminate the Arian party, whose teachings continued to spread. This inaugurated a period of theological strife between Arius and Athanasius. The situation was aggravated by the infiltration of Arianism into the imperial court and its increasing popularity among the populace, whose thinking was more amenable to the simple and pragmatic ideas of Arius. In
addition, Arius expressed his ideas in a series of popular poetic hymns called Thalia (banquet), setting them to music adapted from old, familiar tunes of the ancient Egyptians. These could be heard in the shipyards and all over Alexandria.
Emperor CONSTANTINE I, eager to preserve the unity of his empire, first accepted the verdict of Nicaea, but later wavered in his judgment. He was probably influenced by his Arian sister Constantia and Eusebius of Nicomedia, as well as by the expanding number of Arian followers. At this point Arius seemed to soften his attitude toward the Nicene decision, and the emperor consequently wanted Athanasius to be reconciled with his enemies and to reinstate Arius in church communion. A synod of 335 formally confirmed the reconciliation movement, but Arius died mysteriously in the following year, while the suspicious Athanasius continued to refuse a dubious reconciliation. In the meantime, in 335 the emperor commanded Athanasius to go into exile at Trier in Germany. This proved to be only the first of a series of five exiles of this staunch archbishop, who stood fast by his theology against a movement that survived its author and kept expanding.
The Five Exiles
Athanasius remained in Trier a little more than two years, a period during which he must have composed and developed some of his theological works. With the death of Constantine I in 337, Athanasius and some of the banished Nicene bishops were free to return to their dioceses. Though the people of Alexandria hailed him, Athanasius' return was beset by intrigues from outside. Eusebius of Nicomedia was moved to Constantinople where, as a staunch supporter of Arianism, he had direct access to the imperial court and could influence the emperor against Athanasius. The Arians hoped to depose Athanasius, and in 340, they installed Gregory of Cappadocia in the archiepiscopal throne of Alexandria. Athanasius decided to go into hiding, while his new antagonist reveled in orgies and committed heinous crimes such as causing Philagrius to scourge thirty-four women at church and masterminding the incarceration of other pious Christians. In the midst of such Arian atrocities, Athanasius decided to flee at Easter of the year 340 from Alexandria to Rome.
Thus began his second exile, which lasted three years at the curia of Pope Julius I (337-352) in Rome. Apparently this time the exile was a voluntary one. He was accompanied by a number of Egyptian monks, and perhaps the most significant outcome of that exile was the introduction of the monastic system, which had originated in the Egyptian deserts, to the western Latins. The acceptance of the Egyptian monastic order by the Roman papacy must be regarded as a vital step in the development of Christianity in Europe and the preservation of the Roman heritage in the Middle Ages. On this occasion, too, Athanasius established friendly relations with the Roman see, which recognized his position as archbishop and offered him support throughout his reign. Finally, through the influence of Constantius II (337-361), he was restored to his diocese in Alexandria, now vacated from Arian vestiges by the murder of its Arian occupant, Gregory of Cappadocia, in February 345.
Athanasius' return proved to be an honorable one. He was given imperial letters of congratulations at Aquileia, from where he started the long journey home via Constance, Trier, with its memories of the first exile, and Rome, where Pope Julius offered him an eloquent letter of support. He passed through Hadrianopolis on his journey to Antioch, where he had another opportunity to see the eastern emperor Constantius, who received him with honor. Athanasius refrained from vilifying his opponents to the emperor, but took leave to confront his detractors carrying Constantius' letters to the clergy of Egypt pleading for a unified church. From Antioch Athanasius went to Jerusalem, where he attended a council held in his honor. Athanasius arrived in Alexandria on 24 Babah. Gregory of Nazianzus described the tumultuous welcome by the people, who streamed forth "like the river Nile." Even the dissenting Arian element of the population seemed, for the time being, to have faced the prelate's restoration with charitable clemency.
After these festivities, peace appeared to have reigned for a few years, and as many as 400 bishops, including those substituting for Arians, showered letters of support on Athanasius. Emperor Valens (364-378), it would seem, anathematized Arianism. Nevertheless, it would also seem that the Arians were only procrastinating in the face of irresistible support for Athanasius. Although the western emperor Constans (337-350) lent his support to the great prelate, the eastern emperor Constantius II, being greatly influenced by Eusebius in particular and the Arians in general, turned against him. Constantius condoned and even incited the persecution of Athanasius.
As the restoration festivities began to calm down, the pro-Arian military power started to maneuver against the reestablished orthodoxy. One story is told of a certain General Sebastian, a Manichaean, who came with a batallion of 3,000 soldiers to a cemetery where a company of virgins remained in prayer after the rest of the congregation had dispersed. Sebastian asked the virgins to embrace Arianism. When they refused, he had them stripped and they were thrashed so heavily that some of them died. Such incidents were reputed to have occurred around the churches of Alexandria where the Orthodox bishops were relentlessly pursued, and dozens of them fled. Athanasius stood firm against his opponents, who were led by a certain George the Cappadocian, who intended to replace him on the throne of Alexandria. Athanasius resisted until he could place his case before Constantius, but to his disappointment the emperor issued an imperial missive in which he described the orthodox prelate as a criminal fugitive. Constantius also advised the Ethiopian sovereign to send Frumentius (see ETHIOPIAN PRELATES) founder of the Ethiopian church, back to receive his new ordination, not from Athanasius, but from the Cappodocian George. Frumentius had been consecrated as bishop of Axum by Athanasius prior to 368. Cornered, Athanasius exiled himself from Alexandria again. He joined the increasing number of Coptic monks in the desert. During this exile, which lasted more than six years, he wrote most of his theological works. While keeping contact with his Alexandrian flock through letters of encouragement, he moved from the Nitrian Desert to the Thebaid and lived for a while in the Eastern Desert. He must have spent some time with Saint ANTONY THE GREAT before his death. And it was then that he was able to compose his classic work, the Life of Saint Antony.
With the death of Constantius II in 361, JULIAN THE APOSTATE (361-363) acceded to the imperial throne. Julian had long been contemptuous of the arguments of the Christians, whether orthodox or Arian. The immediate result of Julian's accession was the emergence of the pagan population, who were determined to avenge themselves on George the Cappadocian, who had been determined to exterminate them. The Arian bishop was murdered, and his body was then circulated through the city on a camel. Finally, he was cremated and his remains were thrown into the sea. Although Julian did not favor this gesture at the beginning of his reign, he issued an edict allowing all fugitive bishops to return peacefully to their dioceses. Seizing this opportunity, Athanasius returned to Alexandria on 22 February 362, where he was again met with tumultuous glee by his Orthodox followers.
On his return, Athanasius held a council to resolve all outstanding problems, whether in Alexandria or in Antioch. One of the decisions of that council provided that all who had forfeited their communion with the church could regain it by simply declaring their allegiance to the terms of Nicaea. Those who spoke of three hypostases were found to mean three persons, whereas the Nicene formula prescribed one HYPOSTASIS, the actual Incarnation of the Logos, or assumption of manhood by the Son. Athanasius issued a synodal epistle or tome to the Antiochians about the findings of the council in the hope of achieving church unity. He was unsuccessful because Paulinus, a dissenter, had already been elevated to the episcopate of Antioch, thus starting a schism. In the midst of these internal difficulties, Julian the Apostate denounced all Christian teachings, as well as the right of Athanasius to his episcopal throne. Julian issued a special edict for the expulsion of Athanasius, which was communicated to him by Pythiodorus, a pagan philosopher, on 23 October 362. Thus began the fourth exile of the great prelate. Athanasius left for Memphis and the Thebaid in the year 363. After Julian's death in 363 Athanasius
returned to his episcopal throne.
After his arrival at Alexandria via Hermopolis, where he was hailed by throngs of monks, he received an encouraging letter from the new emperor Jovian (363-364) instructing him to exercise the duties of his episcopal office and prepare a formal statement delineating the Orthodox elements of the faith. Athanasius at once summoned a council, which, under his leadership, framed a synodal epistle that affirmed the Nicene Creed and condemned Arianism and Semi-Arianism, while it denounced the triple definition of the hypostasis and maintained the coequality of the Holy Spirit with the Father and the Son—positions that anticipated the terms of the creed of Constantinople (381).
Armed with these decisions, Athanasius sailed to Antioch, where he was enthusiastically received and where his principles were accepted. The church was consequently united, and even in the West, Pope Liberius (352-355 and 365-366) is known to have made a full declaration of orthodoxy in Rome. In 364, after writing another festal letter at Antioch, Athanasius safely returned to Alexandria shortly before Jovian's death. He was succeeded by Valens (364-378), who was confided the administration of the eastern empire by his brother Valentinian II (375-392). In 365, Valens ordered the expulsion of the bishops that had been allowed to return by Julian. The newly adopted Arian policy caused trouble for the orthodox population and in particular for Athanasius, who stood on the defensive, while the prefect of Alexandria mustered his forces to act against the prelate. Athanasius quietly made his escape through the Church of Saint Dionysius and took refuge for the next four months in a house outside the city. This short period might be considered a fifth and self-inflicted exile. It was terminated by the advent of an imperial notary named Barasides, who came forth with another imperial order for the release of the prelate and his return to his episcopal throne.
From the time of his return to Alexandria until his death in May 373, Athanasius was occupied by disputations against the Arians, by the building of new churches, and in writing some of his final works. He occupied the See of Saint Mark for a total of forty-six years during which he was subjected to persecution that bordered on martyrdom, but his faith in the Nicene Creed was never shaken. According to the Copto-Arabic Synaxarion, his commemoration date occurs on 30 Tut.
In his later years, Athanasius completed his triumph over Arianism, whose exponents were silenced in the Byzantine empire. With the extermination of their teachings from the empire, the splinter of their remaining representatives crossed the Byzantine borders to the realm of the barbarians where they could preach their Arian doctrines in peace to the Goths. Their apostle was ULPHILAS (c. 311-383), originally a Cappadocian, who was consecrated bishop by Eusebius, the Arian bishop of Nicomedia. The Goths remained faithful to Arian precepts even until they descended on the western Roman empire. Their conversion to orthodoxy was a lengthy process in subsequent centuries.
Together with the discomfiture of Arianism and the firm establishment of the Nicene Creed, Athanasius, through his relations with the Pachomian monks and Serapion, was able to give monasticism and ascetic life in Egypt tremendous encouragement and support. Moreover, he was directly responsible for the introduction of monastic rule in the West. As a biographer of Saint Antony, he dedicated his life of the great saint to the people of Gaul and Italy. His theology remained the solid rock on which future generations of theologians continued to build. He was canonized, and the next generation described him as the Apostolic and the Great.
Athanasius is known to have written most of his works in Greek and has been described in A Select Library of the Nicene and Post- Nicene Fathers (1953, p. ixvi) as a Greek father. In fact, the Greek fathers did not know Coptic, and Athanasius, like many educated Copts, was proficient in both Greek and Coptic. Antony and Athanasius must have communicated in Coptic, for Antony did not know Greek.
While Athanasius was still in his twenties, around the year 318, he wrote two short treatises: Against the Gentiles and De Incarnatione Verbi Dei, which became an authoritative theological classic. His thesis in the latter treatise is that by the union of God the Logos with manhood in the person of the Son, Jesus restored to fallen humanity the image of God in which it had been created (Gn. 1:27). By his death and resurrection, Jesus overcame death, which was the consequence of sin. Both treatises predated the outbreak of the Arian controversy in 319. Most of his subsequent work concentrated on the opposition to Arianism beginning with Nicaea.
It is not easy to present a complete bibliography of Athanasius, which has been progressively enriched by new discoveries. Attempts at a compilation of his works have been made by scholars since 1482 when, for the first time, a Latin version of some of his works appeared. Subsequently, two of his genuine works together with a group of spurious ones appeared in Paris in 1520. While rejecting the authenticity of the letters to Serapion, Erasmus edited another collection in 1527; an edition combining the collections of 1520 and 1527 appeared at Lyons in 1532. A more developed Latin edition of all his known works was published by Nannius in 1556, while the first Greek edition by P. Felckmann appeared at Heidelberg in 1608-1612. The Greek text with a Latin version published in Paris in 1627 seems to have superseded all others and may have been supplemented by the one dated 1681 in Leipzig. However, all were overshadowed by the Benedictine edition of 1698 to which B. Montfaucon juxtaposed the Life of the saint. Additional remnants by Montfaucon were compiled in 1707 within the series known as Nova patrum et scriptorum Graecorum collectio. Athanasius' work on the Psalms was edited by N. Antonelli at Rome in 1746 and republished in four folio volumes, which incorporated most of his previous works. Published in English at Oxford in 1842-1844 are the Historical Tracts of St. Athanasius as well as two volumes of Treatises in Controversy with the Arians. His works include the festal letters; his encyclicals; and his special letters to the monks, to Serapion, to the Egyptians and the Libyans, as well as: Apology to Constantius, Apology for His Flight, Apology Against the Arians, History of the Arians, Against the Gentiles, On the Incarnation, Orations and Discourses Against the Arians, Exposition of the Psalms, and Life of Saint Antony.
AZIZ S. ATIYA
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“Honor is a harder master than the law”
They thought Sam Clemens (Mark Twain) was too old to pay everybody back
Sam Clemens, true to his word.
Sam Clemens left his family’s Paris residence, went to New York, and on Wednesday afternoon, April 18,1894 filed for bankruptcy. “Cheer up,” he wrote his wife Livy, “the worst is yet to come.”
Although his travel books and novels were immensely popular, and he had published the bestselling memoirs of Ulysses S. Grant, Clemens (which is how one refers to him about his personal life) was a poor businessman. Moreover, he invested such profits as he made in inventions which didn’t pan out. The worst was a typesetter invented by James W. Paige, an example of which can be seen in the Mark Twain Home, Hartford, Connecticut. If this typesetter had worked, it would have been a wonder, but it was so complicated that it could never be kept working for long. What hopes Clemens had of averting bankruptcy were dashed during the Panic of 1893 when some 16,000 businesses failed.
Clemens owed 96 creditors $94,000, plus $65,000 of his wife’s money which had been loaned to these ventures. This was an overwhelming amount of debt.
At the time of his bankruptcy, Clemens was 58 and in ill health, so his prospects for paying the debts looked bleak. He was fortunate that a friend, Henry Rogers, a partner in Standard Oil, offered to help deal with creditors and handle his finances. One of Rogers’ achievements was to keep the literary copyrights out of creditor negotiations and in Livy’s name. The Clemens and Rogers agreed that the creditors must be repaid in full. As Clemens explained, “honor is a harder master than the law.”
By February 1895, Clemens finished two books: Tom Sawyer, Detective and Joan of Arc. He told Rogers, “I am in a sort of physical collapse.”
Rogers negotiated the publishing contracts, but it soon became clear that publishing income alone wouldn’t be enough to keep creditors at bay. Clemens had to go back on the lecture circuit which he had done successfully before. A world tour was the best bet to pay off his creditors, but such tours were exhausting. Travel was slow. There were many connections to make, and every connection meant hauling lots of luggage, especially for women – because of Sam’s poor health, it made sense for him to have Livy and daughter Clara along.
“Most grime on clothing was external, caused by street filth and air pollution,” explained scholar Robert Cooper who chronicled the itinerary of the world tour in Around the World with Mark Twain (2000). Horse-drawn vehicles splashed clothes with mud. Horse droppings, which were hard to avoid, and the dirt and mud of the streets, then mostly unpaved, soiled trouser cuffs or the hems of long dresses. Coal smoke, from cooking and heating as well as from industrial users, smudged clothes as well. So the Victorian custom of changing for the evening meal was not an affectation…
“When washing or cleaning a dress, it was common practice to open up the seams, even elaborate garments with forty or fifty yards of fabric, lay out all the sections on a clean table, smooth out the creases, scrub each side with soapy water and then rinse (or, if the dress could not be washed, apply a cleaning agent such as gin), before sewing the sections back together. Since cleaning was often time-consuming, travelers had to pack enough garments so that they could wear clean clothes while dirt was being removed from the soiled ones.” As a result, a tour such as Clemens contemplated meant going around the world with 16 pieces of hand luggage and many large steamer trunks.
His daughter Clara later recalled how he anguished about “the hellish struggle it was to settle on making that lecture trip around the world? How we fought the idea, the horrible idea, the heart-breaking idea…I am almost an old man, with ill health, carbuncles, bronchitis and rheumatism…with patience worn to rags, I was to pack my bag and be jolted around the devil’s universe…”
Clemens wrote his friend James B. Pond, one of America’s most successful lecture who had arranged a previous tour for him in 1884 and 1885. Pond agreed to range another one across the northern United States and to accompany him, so he could concentrate on his presentations. Pond would get a quarter of the profits. As biographer Justin Kaplan reported in Mr. Clemens and Mark Twain (1966), “confined to his bed for three weeks with carbuncles and gout and half expecting to start his travels on a stretcher, Clemens prepared his programs and worked over his itinerary.” He said, “I am confident that if I live I can pay off the last debt within four years.
What would the program be? A lecture? A reading? Earlier in his career, he recalled in his autobiography (Charles Neider edition), “I supposed it would only be necessary to do like Dickens – get out on the platform and read from the book. I did that and made a botch of it. Written things are not for speech; their form is literary; they are stiff, inflexible, and will not lend themselves to happy and effective delivery with the tongue – where their purpose is to merely entertain, not instruct; they have to be limbered up, broken up, colloquialized and turned into the common forms of unpremeditated talk – otherwise they will bore the house, not entertain it…I memorized those pieces, and in delivering them from the platform they soon transformed themselves into flexible talk, with all their obstructing preciseness and formalities gone out of them for good.”
He added, “in reading from the book you are telling another person’s tale at secondhand; you are a mimic and not the person involved; you are an artificiality, not a reality; whereas in telling the tale without the book, you absorb the character and presently become the man himself, just as is the case with the actor.”
Fred W. Lorch, whose book The Trouble Begins at Eight (1968) chronicled Mark Twain’s lecture career, reported that he knew “people wanted a lecturer to give them something solid, something educational, something that would improve them. And that was what he had come to do. He proposed to teach morals by the use of illustrations. He had a theory that a person should prize as priceless every crime, every transgression he commits – that is, the lesson he derives from it. By impressing the lesson of the crime upon his mind and heart, a person would never commit that crime again…
“Mark Twain then proceeded to illustrate some of the principles which he asserted he had learned in his own rise toward moral perfection (he was now, he said, more than two-thirds on his way up there and had not much further to go) by introducing various stories and anecdotes from his own writings.” Since he would be scheduled to lecture two or three times in some cities and towns, he needed an adequate repertoire of stories.
“The Jumping Frog and the story of the Mexican Plug,” Lorch continued, “were offered to teach a person never to put faith in a passing stranger; Jim Baker’s Blue Jay, whatever you do, do with all your heart; The History of a Campaign Trail That Failed, discretion is the better part of valor; Tom Sawyer’s Crusade To Rescue the Holy Land, don’t argue matters beyond your comprehension; the Corpse in His Father’s Office, learn to gauge you courage early; The Awful German Language, the necessity of teaching patience; Huck Helps Jim Escape, a sound heart is better than a deformed conscience; The Christening of Mary Ann, don’t jump to conclusions; The Smallpox story, a fellow has to start early in life if he wants to do right…”
On July 14, 1895, Sam, Livy and Clara Clemens boarded a train from Elmira, New York, where the family spent summers, to Cleveland where the tour would begin the next day. Then they went to Sault Ste. Marie, Mackinac and Petoskey, Michigan; Deluth, Minneapolis and St. Paul, Minnesota; Winnipeg, Manitoba; Butte, Anaconda, Helena, Missoula and Great Falls, Montana; Spokane, Olympia, Tacoma and Seattle, Washington; Portland, Oregon; and Vancouver and Victoria, British Columbia. Ill health plagued Clemens much of the way, and he missed many dinners which had been planned in his honor. But he played to packed houses, sometimes over a thousand people at a time, and by the time he was done with the U.S. leg of his tour he was able to remit some $5,000 to his creditors.
People welcomed him everywhere, which helped lift his spirits amidst the crisis of bankruptcy. Everybody knew his books Innocents Abroad (1869), The Adventures of Tom Sawyer (1876) and The Adventures of Huckleberry Finn (1884). And of course, Mark Twain had become a famous attraction on the lecture circuit ever since his first lecture in San Francisco back on October 2, 1866. As he told a newspaper reporter at the end of his U.S. tour, “Lecturing is gymnastics, chest-expander, medicine, mind healer, blues destroyer, all in one.”
Mark Twain on stage was quite a sight. According to Lorch, “For a man sixty years of age and recently ill, his appearance was truly impressive. Bushy white hair circled his head like a halo. The brows were shaggy and thick, the mustache dropped, strongly aquiline, the thin delicate nostrils suggesting nervous sensibility. The eyes, wonderfully keen and piercing, looked out gently from a furrowed, intellectual face. He appeared in full evening dress with a wide unornamental expanse of white shirt, and it occurred to some in the audience that he looked far too civilized to be the author of Roughing It.
“Emerging from the wings, his left hand thrust deep into his trouser pocket, he sauntered slowly toward the reading stand like a man out for a stroll who presently looks up and finds he has company. At times he walked out on stage carrying a glass of water in such a casual and unstudied way as to appear that it was a perfectly natural thing for him to do. Then, stepping to the side of the stand, he stood with folded arms gazing at the audience with quizzical look that betrayed no hint of a smile, and waited for the audience to quiet down…
“His speech was so distinct, clear, and penetrating that he could be heard everywhere in the hall, except occasionally when he lowered his voice. What impressed his audiences most, however, about his manner of speech was his slow, leisurely drawl (a Yankee drawl, as reporters characterized it) and his measured manner of speaking, never repeating or withdrawing a word once uttered.”
From Victoria, British Columbia, the Clemens sailed for Australia. They stopped at Hawaii but couldn’t get off the ship because of a cholera outbreak in Honolulu. Carlyle G. Smythe, one of Australia’s best-known theatrical managers, had arranged a lecture tour through the British Commonwealth. The Clemens arrived in Sydney, Australia on Monday, September 16, 1895. He performed there, in Melbourne, Adelaide, Horsham, Stawell, Ballarat, Bendigo, Maryborough, Geelong and Prahan. On to New Zealand where he performed in Invercargill, Dunedin, Timaru, Oamaru, Christchurch, Auckland and Wellington. Then back to Melbourne and Sydney for more “At Homes.” According to biographer Kaplan, the first two weeks in Australia netted $2,200 for creditors.
“Foreign audiences, like those at home,” Lorch wrote, “were greatly impressed by the distinguished humorist’s subtle and masterful resources for capturing and holding audience interest. They noted, for example, the skillfully used pause, after which invariably followed a graphic word or a meaningful phrase purposely withheld; the serious, imperturbable demeanor; the dead pan expression; the occasional groping for a word, though no one imagined that the word was not instantly at his command; and finally the complete naturalness and simplicity of manner by which he concealed the devices which made them laugh.
“Try as they might, reporters confessed they found it almost impossible adequately to describe the effect which Mark Twain’s performance exerted upon them. What was the magic, they asked themselves, which enabled him for nearly two hours to keep them under his sway? As a lecture it was so utterly out of the ordinary, so unique in material, so delightfully rambling and inconsequential in treatment that no standard came to mind with which it could be critically compared. Even if one were to reproduce his talk verbatim, the magic would be lost; for the reporters quickly sensed that the charm was not primarily in the stories themselves, but rather in the manner of telling, in the techniques and in the whole speaking personality of the humorist. All these blended themselves together so simply and naturally and effortlessly that no amount of specification could supply a true and adequate notion of the effect.
“That Mark Twain was not merely a funny man, a mere laughter-maker, they also perceived. This man sensed the tears in human affairs. He perceived the acute suffering that afflicted the souls of men."
Especially on his overseas tour, Clemens insisted on having time between performances when he could go sight-seeing, because he intended to write a book about his experiences, expected to bring in some money. His book Innocents Abroad, about a tour he took to the Holy Land, had been his first bestseller. Roughing It (1872), another chronicle of his travels, had done well.
At every stop, Clemens met reporters, and many of the details of his world tour came from local newspaper accounts which Fred Lorch pieced together. “The reporters,” he wrote, “frequently observed that he spent a great deal of time in collecting information about the localities he visited and in writing and assembling his notes.”
He found it tough to work because so many people wanted to meet him. Lorch reported, “American consular agents, prominent local and state officials, high ranking military personnel, clubs and organizations which had sponsored his lectures – all invited him to dinners, receptions, and other affairs, often late in the evening after the lecture. If he accepted, it usually meant exhausting hours that robbed him of rest and drained his energies.”
From Australia, the Clemens sailed to Colombo (Ceylon), then India. He performed in Bombay, Poona, Baroda, Allahabad, Calcutta, Darjeeling, Muzaffapur, Lucknow, Cawnpore, Agra, Lahore and Rawalpindi. On to South Africa and performances in Durban, Pietermaritzburg, Johannesburg, Pretoria, Krugersdorp, Bloemfontein, Queenstown, King William’s Town, East London, Port Elizabeth, Grahamstown, Kimberly and Cape Town.
The Clemens sailed for Southampton, England on July 15, 1896. Three weeks after they arrived, there was terrible news: daughter Susy, who with her sister Jean had remained in the states with their aunt, had died of meningitis at their Hartford, Connecticut home. She was just 24. Clemens didn’t have the heart to give any more lectures. “It is one of the mysteries of our nature,” he said, “that a man, all unprepared, can receive a thunder-stroke like that and live.”
He poured his energies into Following the Equator which he finished in May 1897. More than 30,000 copies were sold, and Clemens remarked, “Land, we are glad to see those debts diminishing. For the first time in my life I am getting more pleasure from paying money out than from pulling it in.” Clemens repeatedly thanked Rogers for his help.
By January 1898, Rogers cabled Clemens: “The creditors have all been paid a hundred cents on the dollar.” So Clemens had stepped up to take responsibility for his financial crisis and come through it with his integrity and peace of mind.
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We have within the modern period very many examples which enable us to study the evolution of legendary miracles. Out of these I will select but one, which is chosen because it is the life of one of the most noble and devoted men in the history of humanity, one whose biography is before the world with its most minute details - in his own letters, in the letters of his associates, in contemporary histories, and in a multitude of biographies: this man is St. Francis Xavier. From these sources I draw the facts now to be given, but none of them are of Protestant origin; every source from which I shall draw is Catholic and Roman, and published under the sanction of the Church.
Born a Spanish noble, Xavier at an early age cast aside all ordinary aims, devoted himself to study, was rapidly advanced to a professorship at Paris, and in this position was rapidly winning a commanding influence, when he came under the sway of another Spaniard even greater, though less brilliantly endowed, than himself - Ignatius Loyola, founder of the Society of Jesus. The result was that the young professor sacrificed the brilliant career on which he had entered at the French capital, went to the far East as a simple missionary, and there devoted his remaining years to redeeming the lowest and most wretched of our race.
Among the various tribes, first in lower India and afterward in Japan, he wrought untiringly - toiling through village after village, collecting the natives by the sound of a hand-bell, trying to teach them the simplest Christian formulas; and thus he brought myriads of them to a nominal Confession of the Christian faith. After twelve years of such efforts, seeking new conquests for religion, he sacrificed his life on the desert island of San Chan.
During his career as a missionary he wrote great numbers of letters, which were preserved and have since been published; and these, with the letters of his contemporaries, exhibit clearly all the features of his life. His own writings are very minute, and enable us to follow him fully. No account of a miracle wrought by him appears either in his own letters or in any contemporary document. At the outside, but two or three things occurred in his whole life, as exhibited so fully by himself and his contemporaries, for which the most earnest devotee could claim anything like Divine interposition; and these are such as may be read in the letters of very many fervent missionaries, Protestant as well as Catholic. For example, in the beginning of his career, during a journey in Europe with an ambassador, one of the servants in fording a stream got into deep water and was in danger of drowning. Xavier tells us that the ambassador prayed very earnestly, and that the man finally struggled out of the stream. But within sixty years after his death, at his canonization, and by various biographers, this had been magnified into a miracle, and appears in the various histories dressed out in glowing colours. Xavier tells us that the ambassador prayed for the safety of the young man; but his biographers tell us that it was Xavier who prayed, and finally, by the later writers, Xavier is represented as lifting horse and rider out of the stream by a clearly supernatural act.
Still another claim to miracle is based upon his arriving at Lisbon and finding his great colleague, Simon Rodriguez, ill of fever. Xavier informs us in a very simple way that Rodriguez was so overjoyed to see him that the fever did not return. This is entirely similar to the cure which Martin Luther wrought upon Melanchthon. Melanchthon had broken down and was supposed to be dying, when his joy at the long-delayed visit of Luther brought him to his feet again, after which he lived for many years.
Again, it is related that Xavier, finding a poor native woman very ill, baptized her, saying over her the prayers of the Church, and she recovered.
Two or three occurrences like these form the whole basis for the miraculous account, so far as Xavier's own writings are concerned.
Of miracles in the ordinary sense of the word there is in these letters of his no mention. Though he writes of his doings with especial detail, taking evident pains to note everything which he thought a sign of Divine encouragement, he says nothing of his performing miracles, and evidently knows nothing of them. This is clearly not due to his unwillingness to make known any token of Divine favour. As we have seen, he is very prompt to report anything which may be considered an answer to prayer or an evidence of the power of religious means to improve the bodily or spiritual health of those to whom he was sent.
Nor do the letters of his associates show knowledge of any miracles wrought by him. His brother missionaries, who were in constant and loyal fellowship with him, make no allusions to them in their communications with each other or with their brethren in Europe.
Of this fact we have many striking evidences. Various collections of letters from the Jesuit missionaries in India and the East generally, during the years of Xavier's activity, were published, and in not one of these letters written during Xavier's lifetime appears any account of a miracle wrought by him. As typical of these collections we may take perhaps the most noted of all, that which was published about twenty years after Xavier's death by a Jesuit father, Emanuel Acosta.
The letters given in it were written by Xavier and his associates not only from Goa, which was the focus of all missionary effort and the centre of all knowledge regarding their work in the East, but from all other important points in the great field. The first of them were written during the saint's lifetime, but, though filled with every sort of detail regarding missionary life and work, they say nothing regarding any miracles by Xavier.
The same is true of various other similar collections published during the sixteenth and seventeenth centuries. In not one of them does any mention of a miracle by Xavier appear in a letter from India or the East contemporary with him.
This silence regarding his miracles was clearly not due to any ``evil heart of unbelief.'' On the contrary, these good missionary fathers were prompt to record the slightest occurrence which they thought evidence of the Divine favour: it is indeed touching to see how eagerly they grasp at the most trivial things which could be thus construed.
Their ample faith was fully shown. One of them, in Acosta's collection, sends a report that an illuminated cross had been recently seen in the heavens; another, that devils had been cast out of the natives by the use of holy water; another, that various cases of disease had been helped and even healed by baptism; and sundry others sent reports that the blind and dumb had been restored, and that even lepers had been cleansed by the proper use of the rites of the Church; but to Xavier no miracles are imputed by his associates during his life or during several years after his death.
On the contrary, we find his own statements as to his personal limitations, and the difficulties arising from them, fully confirmed by his brother workers. It is interesting, for example, in view of the claim afterward made that the saint was divinely endowed for his mission with the ``gift of tongues,'' to note in these letters confirmation of Xavier's own statement utterly disproving the existence of any such Divine gift, and detailing the difficulties which he encountered from his want of knowing various languages, and the hard labour which he underwent in learning the elements of the Japanese tongue.
Until about ten years after Xavier's death, then, as Emanuel Acosta's publication shows, the letters of the missionaries continued without any indication of miracles performed by the saint. Though, as we shall see presently, abundant legends had already begun to grow elsewhere, not one word regarding these miracles came as yet from the country which, according to later accounts accepted and sanctioned by the Church, was at this very period filled with miracles; not the slightest indication of them from the men who were supposed to be in the very thick of these miraculous manifestations.
But this negative evidence is by no means all. There is also positive evidence - direct testimony from the Jesuit order itself - that Xavier wrought no miracles.
For not only did neither Xavier nor his co-workers know anything of the mighty works afterward attributed to him, but the highest contemporary authority on the whole subject, a man in the closest correspondence with those who knew most about the saint, a member of the Society of Jesus in the highest standing and one of its accepted historians, not only expressly tells us that Xavier wrought no miracles, but gives the reasons why he wrought none.
This man was Joseph Acosta, a provincial of the Jesuit order, its visitor in Aragon, superior at Valladolid, and finally rector of the University of Salamanca. In 1571, nineteen years after Xavier's death, Acosta devoted himself to writing a work mainly concerning the conversion of the Indies, and in this he refers especially and with the greatest reverence to Xavier, holding him up as an ideal and his work as an example.
But on the same page with this tribute to the great missionary Acosta goes on to discuss the reasons why progress in the world's conversion is not so rapid as in the early apostolic times, and says that an especial cause why apostolic preaching could no longer produce apostolic results ``lies in the missionaries themselves, because there is now no power of working miracles.'' He then asks, ``Why should our age be so completely destitute of them?'' This question he answers at great length, and one of his main contentions is that in early apostolic times illiterate men had to convert the learned of the world, whereas in modern times the case is reversed, learned men being sent to convert the illiterate; and hence that ``in the early times miracles were necessary, but in our time they are not.''
This statement and argument refer, as we have seen, directly to Xavier by name, and to the period covered by his activity and that of the other great missionaries of his time. That the Jesuit order and the Church at large thought this work of Acosta trustworthy is proved by the fact that it was published at Salamanca a few years after it was written, and republished afterward with ecclesiastical sanction in France. Nothing shows better than the sequel how completely the evolution of miraculous accounts depends upon the intellectual atmosphere of any land and time, and how independent it is of fact.
For, shortly after Xavier's heroic and beautiful death in 1552, stories of miracles wrought by him began to appear. At first they were few and feeble; and two years later Melchior Nunez, Provincial of the Jesuits in the Portuguese dominions, with all the means at his command, and a correspondence extending throughout Eastern Asia, had been able to hear of but three. These were entirely from hearsay. First, John Deyro said he knew that Xavier had the gift of prophecy; but, unfortunately, Xavier himself had reprimanded and cast off Deyro for untruthfulness and cheatery. Secondly, it was reported vaguely that at Cape Comorin many persons affirmed that Xavier had raised a man from the dead. Thirdly, Father Pablo de Santa Fe had heard that in Japan Xavier had restored sight to a blind man. This seems a feeble beginning, but little by little the stories grew, and in 1555 De Quadros, Provincial of the Jesuits in Ethiopia, had heard of nine miracles, and asserted that Xavier had healed the sick and cast out devils. The next year, being four years after Xavier's death, King John III of Portugal, a very devout man, directed his viceroy Barreto to draw up and transmit to him an authentic account of Xavier's miracles, urging him especially to do the work ``with zeal and speedily.'' We can well imagine what treasures of grace an obsequious viceroy, only too anxious to please a devout king, could bring together by means of the hearsay of ignorant, compliant natives through all the little towns of Portuguese India.
But the letters of the missionaries who had been co-workers or immediate successors of Xavier in his Eastern field were still silent as regards any miracles by him, and they remained silent for nearly ten years. In the collection of letters published by Emanuel Acosta and others no hint at any miracles by him is given, until at last, in 1562, fully ten years after Xavier's death, the first faint beginnings of these legends appear in them.
At that time the Jesuit Almeida, writing at great length to the brethren, stated that he had found a pious woman who believed that a book left behind by Xavier had healed sick folk when it was laid upon them, and that he had met an old man who preserved a whip left by the saint which, when properly applied to the sick, had been found good both for their bodies and their souls. From these and other small beginnings grew, always luxuriant and sometimes beautiful, the vast mass of legends which we shall see hereafter.
This growth was affectionately garnered by the more zealous and less critical brethren in Europe until it had become enormous; but it appears to have been thought of little value by those best able to judge.
For when, in 1562, Julius Gabriel Eugubinus delivered a solemn oration on the condition and glory of the Church, before the papal legates and other fathers assembled at the Council of Trent, while he alluded to a multitude of things showing the Divine favour, there was not the remotest allusion to the vast multitude of miracles which, according to the legends, had been so profusely lavished on the faithful during many years, and which, if they had actually occurred, formed an argument of prodigious value in behalf of the special claims of the Church.
The same complete absence of knowledge of any such favours vouchsafed to the Church, or at least of any belief in them, appears in that great Council of Trent among the fathers themselves. Certainly there, if anywhere, one might on the Roman theory expect Divine illumination in a matter of this kind. The presence of the Holy Spirit in the midst of it was especially claimed, and yet its members, with all their spiritual as well as material advantages for knowing what had been going on in the Church during the previous thirty years, and with Xavier's own friend and colleague, Laynez, present to inform them, show not the slightest sign of any suspicion of Xavier's miracles. We have the letters of Julius Gabriel to the foremost of these fathers assembled at Trent, from 1557 onward for a considerable time, and we have also a multitude of letters written from the Council by bishops, cardinals, and even by the Pope himself, discussing all sorts of Church affairs, and in not one of these is there evidence of the remotest suspicion that any of these reports, which they must have heard, regarding Xavier's miracles, were worthy of mention.
Here, too, comes additional supplementary testimony of much significance. With these orations and letters, Eugubinus gives a Latin translation of a letter, ``on religious affairs in the Indies,'' written by a Jesuit father twenty years after Xavier's death. Though the letter came from a field very distant from that in which Xavier laboured, it was sure, among the general tokens of Divine favour to the Church and to the order, on which it dwelt, to have alluded to miracles wrought by Xavier had there been the slightest ground for believing in them; but no such allusion appears.
So, too, when in 1588, thirty-six years after Xavier's death, the Jesuit father Maffei, who had been especially conversant with Xavier's career in the East, published his History of India, though he gave a biography of Xavier which shows fervent admiration for his subject, he dwelt very lightly on the alleged miracles. But the evolution of miraculous legends still went on. Six years later, in 1594, Father Tursellinus published his Life of Xavier, and in this appears to have made the first large use of the information collected by the Portuguese viceroy and the more zealous brethren. This work shows a vast increase in the number of miracles over those given by all sources together up to that time. Xavier is represented as not only curing the sick, but casting out devils, stilling the tempest, raising the dead, and performing miracles of every sort.
In 1622 came the canonization proceedings at Rome. Among the speeches made in the presence of Pope Gregory XV, supporting the claims of Xavier to saintship, the most important was by Cardinal Monte. In this the orator selects out ten great miracles from those performed by Xavier during his lifetime and describes them minutely. He insists that on a certain occasion Xavier, by the sign of the cross, made sea-water fresh, so that his fellow-passengers and the crew could drink it; that he healed the sick and raised the dead in various places; brought back a lost boat to his ship; was on one occasion lifted from the earth bodily and transfigured before the bystanders; and that, to punish a blaspheming town, he caused an earthquake and buried the offenders in cinders from a volcano: this was afterward still more highly developed, and the saint was represented in engravings as calling down fire from heaven and thus destroying the town.
The most curious miracle of all is the eighth on the cardinal's list. Regarding this he states that, Xavier having during one of his voyages lost overboard a crucifix, it was restored to him after he had reached the shore by a crab.
The cardinal also dwelt on miracles performed by Xavier's relics after his death, the most original being that sundry lamps placed before the image of the saint and filled with holy water burned as if filled with oil.
This latter account appears to have deeply impressed the Pope, for in the Bull of Canonization issued by virtue of his power of teaching the universal Church infallibly in all matters pertaining to faith and morals, His Holiness dwells especially upon the miracle of the lamp filled with holy water and burning before Xavier's image.
Xavier having been made a saint, many other Lives of him appeared, and, as a rule, each surpassed its predecessor in the multitude of miracles. In 1622 appeared that compiled and published under the sanction of Father Vitelleschi, and in it not only are new miracles increased, but some old ones are greatly improved. One example will suffice to show the process. In his edition of 1596, Tursellinus had told how, Xavier one day needing money, and having asked Vellio, one of his friends, to let him have some, Vellio gave him the key of a safe containing thirty thousand gold pieces. Xavier took three hundred and returned the key to Vellio; whereupon Vellio, finding only three hundred pieces gone, reproached Xavier for not taking more, saying that he had expected to give him half of all that the strong box contained. Xavier, touched by this generosity, told Vellio that the time of his death should be made known to him, that he might have opportunity to repent of his sins and prepare for eternity. But twenty-six years later the Life of Xavier published under the sanction of Vitelleschi, giving the story, says that Vellio on opening the safe found that all his money remained as he had left it, and that none at all had disappeared; in fact, that there had been a miraculous restitution. On his blaming Xavier for not taking the money, Xavier declares to Vellio that not only shall he be apprised of the moment of his death, but that the box shall always be full of money. Still later biographers improved the account further, declaring that Xavier promised Vellio that the strong box should always contain money sufficient for all his needs. In that warm and uncritical atmosphere this and other legends grew rapidly, obedient to much the same laws which govern the evolution of fairy tales.
In 1682, one hundred and thirty years after Xavier's death, appeared his biography by Father Bouhours; and this became a classic. In it the old miracles of all kinds were enormously multiplied, and many new ones given. Miracles few and small in Tursellinus became many and great in Bouhours. In Tursellinus, Xavier during his life saves one person from drowning, in Bouhours he saves during his life three; in Tursellinus, Xavier during his life raises four persons from the dead, in Bouhours fourteen; in Tursellinus there is one miraculous supply of water, in Bouhours three; in Tursellinus there is no miraculous draught of fishes, in Bouhours there is one; in Tursellinus, Xavier is transfigured twice, in Bouhours five times: and so through a long series of miracles which, in the earlier lives appearing either not at all or in very moderate form, are greatly increased and enlarged by Tursellinus, and finally enormously amplified and multiplied by Father Bouhours.
And here it must be borne in mind that Bouhours, writing ninety years after Tursellinus, could not have had access to any new sources. Xavier had been dead one hundred and thirty years, and of course all the natives upon whom he had wrought his miracles, and their children and grandchildren, were gone. It can not then be claimed that Bouhours had the advantage of any new witnesses, nor could he have had anything new in the way of contemporary writings; for, as we have seen, the missionaries of Xavier's time wrote nothing regarding his miracles, and certainly the ignorant natives of India and Japan did not commit any account of his miracles to writing. Nevertheless, the miracles of healing given in Bouhours were more numerous and brilliant than ever. But there was far more than this. Although during the lifetime of Xavier there is neither in his own writings nor in any contemporary account any assertion of a resurrection from the dead wrought by him, we find that shortly after his death stories of such resurrections began to appear. A simple statement of the growth of these may throw some light on the evolution of miraculous accounts generally. At first it was affirmed that some people at Cape Comorin said that he had raised one person; then it was said that there were two persons; then in various authors - Emanuel Acosta, in his commentaries written as an afterthought nearly twenty years after Xavier's death, De Quadros, and others - the story wavers between one and two cases; finally, in the time of Tursellinus, four cases had been developed. In 1622, at the canonization proceedings, three were mentioned; but by the time of Father Bouhours there were fourteen - all raised from the dead by Xavier himself during his lifetime - and the name, place, and circumstances are given with much detail in each case.
It seems to have been felt as somewhat strange at first that Xavier had never alluded to any of these wonderful miracles; but ere long a subsidiary legend was developed, to the effect that one of the brethren asked him one day if he had raised the dead, whereat he blushed deeply and cried out against the idea, saying: ``And so I am said to have raised the dead! What a misleading man I am! Some men brought a youth to me just as if he were dead, who, when I commanded him to arise in the name of Christ, straightway arose.''
Noteworthy is the evolution of other miracles. Tursellinus, writing in 1594, tells us that on the voyage from Goa to Malacca, Xavier having left the ship and gone upon an island, was afterward found by the persons sent in search of him so deeply absorbed in prayer as to be unmindful of all things about him. But in the next century Father Bouhours develops the story as follows: ``The servants found the man of God raised from the ground into the air, his eyes fixed upon heaven, and rays of light about his countenance.''
Instructive, also, is a comparison between the successive accounts of his noted miracle among the Badages at Travancore, in 1544 Xavier in his letters makes no reference to anything extraordinary; and Emanuel Acosta, in 1571, declares simply that ``Xavier threw himself into the midst of the Christians, that reverencing him they might spare the rest.'' The inevitable evolution of the miraculous goes on; and twenty years later Tursellinus tells us that, at the onslaught of the Badages, ``they could not endure the majesty of his countenance and the splendour and rays which issued from his eyes, and out of reverence for him they spared the others.'' The process of incubation still goes on during ninety years more, and then comes Father Bouhours's account. Having given Xavier's prayer on the battlefield, Bouhours goes on to say that the saint, crucifix in hand, rushed at the head of the people toward the plain where the enemy was marching, and ``said to them in a threatening voice, `I forbid you in the name of the living God to advance farther, and on His part command you to return in the way you came.' These few words cast a terror into the minds of those soldiers who were at the head of the army; they remained confounded and without motion. They who marched afterward, seeing that the foremost did not advance, asked the reason of it. The answer was returned from the front ranks that they had before their eyes an unknown person habited in black, of more than human stature, of terrible aspect, and darting fire from his eyes.... They were seized with amazement at the sight, and all of them fled in precipitate confusion.''
Curious, too, is the after-growth of the miracle of the crab restoring the crucifix. In its first form Xavier lost the crucifix in the sea, and the earlier biographers dwell on the sorrow which he showed in consequence; but the later historians declare that the saint threw the crucifix into the sea in order to still a tempest, and that, after his safe getting to land, a crab brought it to him on the shore. In this form we find it among illustrations of books of devotion in the next century.
But perhaps the best illustration of this evolution of Xavier's miracles is to be found in the growth of another legend; and it is especially instructive because it grew luxuriantly despite the fact that it was utterly contradicted in all parts of Xavier's writings as well as in the letters of his associates and in the work of the Jesuit father, Joseph Acosta.
Throughout his letters, from first to last, Xavier constantly dwells upon his difficulties with the various languages of the different tribes among whom he went. He tells us how he surmounted these difficulties: sometimes by learning just enough of a language to translate into it some of the main Church formulas; sometimes by getting the help of others to patch together some pious teachings to be learned by rote; sometimes by employing interpreters; and sometimes by a mixture of various dialects, and even by signs. On one occasion he tells us that a very serious difficulty arose, and that his voyage to China was delayed because, among other things, the interpreter he had engaged had failed to meet him.
In various Lives which appeared between the time of his death and his canonization this difficulty is much dwelt upon; but during the canonization proceedings at Rome, in the speeches then made, and finally in the papal bull, great stress was laid upon the fact that Xavier possessed the gift of tongues. It was declared that he spoke to the various tribes with ease in their own languages. This legend of Xavier's miraculous gift of tongues was especially mentioned in the papal bull, and was solemnly given forth by the pontiff as an infallible statement to be believed by the universal Church. Gregory XV having been prevented by death from issuing the Bull of Canonization, it was finally issued by Urban VIII; and there is much food for reflection in the fact that the same Pope who punished Galileo, and was determined that the Inquisition should not allow the world to believe that the earth revolves about the sun, thus solemnly ordered the world, under pain of damnation, to believe in Xavier's miracles, including his ``gift of tongues,'' and the return of the crucifix by the pious crab. But the legend was developed still further: Father Bouhours tells us, ``The holy man spoke very well the language of those barbarians without having learned it, and had no need of an interpreter when he instructed.'' And, finally, in our own time, the Rev. Father Coleridge, speaking of the saint among the natives, says, ``He could speak the language excellently, though he had never learned it.''
In the early biography, Tursellinus writes. ``Nothing was a greater impediment to him than his ignorance of the Japanese tongues; for, ever and anon, when some uncouth expression offended their fastidious and delicate ears, the awkward speech of Francis was a cause of laughter.'' But Father Bouhours, a century later, writing of Xavier at the same period, says, ``He preached in the afternoon to the Japanese in their language, but so naturally and with so much ease that he could not be taken for a foreigner.''
And finally, in 1872, Father Coleridge, of the Society of Jesus, speaking of Xavier at this time, says, ``He spoke freely, flowingly, elegantly, as if he had lived in Japan all his life.''
Nor was even this sufficient: to make the legend complete, it was finally declared that, when Xavier addressed the natives of various tribes, each heard the sermon in his own language in which he was born.
All this, as we have seen, directly contradicts not only the plain statements of Xavier himself, and various incidental testimonies in the letters of his associates, but the explicit declaration of Father Joseph Acosta. The latter historian dwells especially on the labour which Xavier was obliged to bestow on the study of the Japanese and other languages, and says, ``Even if he had been endowed with the apostolic gift of tongues, he could not have spread more widely the glory of Christ.''
It is hardly necessary to attribute to the orators and biographers generally a conscious attempt to deceive. The simple fact is, that as a rule they thought, spoke, and wrote in obedience to the natural laws which govern the luxuriant growth of myth and legend in the warm atmosphere of love and devotion which constantly arises about great religious leaders in times when men have little or no knowledge of natural law, when there is little care for scientific evidence, and when he who believes most is thought most meritorious.
These examples will serve to illustrate the process which in thousands of cases has gone on from the earliest days of the Church until a very recent period. Everywhere miraculous cures became the rule rather than the exception throughout Christendom. |
Active galactic nucleus
An active galactic nucleus (AGN) is a compact region at the centre of a galaxy that has a much higher than normal luminosity over at least some portion, and possibly all, of the electromagnetic spectrum. Such excess emission has been observed in the radio, infrared, optical, ultra-violet, X-ray and gamma ray wavebands. A galaxy hosting an AGN is called an active galaxy. The radiation from AGN is believed to be a result of accretion of mass by a supermassive black hole at the centre of its host galaxy. AGN are the most luminous and persistent sources of electromagnetic radiation in the universe, and as such can be used as a means of discovering distant objects; their evolution as a function of cosmic time also puts constraints on models of the cosmos.
Models of the active nucleus
For a long time it has been argued that an AGN must be powered by accretion of mass onto massive black holes (106 to 1010 times the Solar mass). AGN are both compact and persistently extremely luminous. Accretion can potentially give very efficient conversion of potential and kinetic energy to radiation, and a massive black hole has a high Eddington luminosity, and as a result, it can provide the observed high persistent luminosity. Supermassive black holes are now believed to exist in the centres of most if not all massive galaxies. Evidence for that is that the mass of the black hole correlates well with the velocity dispersion of the galactic bulge (the M-sigma relation) or with bulge luminosity (e.g.). Thus AGN-like characteristics are expected whenever a supply of material for accretion comes within the sphere of influence of the central black hole.
Accretion disc
In the standard model of AGN, cold material close to a black hole forms an accretion disc. Dissipative processes in the accretion disc transport matter inwards and angular momentum outwards, while causing the accretion disc to heat up. The expected spectrum of an accretion disc peaks in the optical-ultraviolet waveband; in addition, a corona of hot material forms above the accretion disc and can inverse-Compton scatter photons up to X-ray energies. The radiation from the accretion disc excites cold atomic material close to the black hole and this in turn radiates at particular emission lines. A large fraction of the AGN's radiation may be obscured by interstellar gas and dust close to the accretion disc, but (in a steady-state situation) this will be re-radiated at some other waveband, most likely the infrared.
Relativistic jets
Some accretion discs produce jets of twin, highly collimated, and fast outflows that emerge in opposite directions from close to the disc. The direction of the jet ejection is determined either by the angular momentum axis of the accretion disc or the spin axis of the black hole. The jet production mechanism and indeed the jet composition on very small scales are not understood at present due to the low resolution of astronomical instruments, and as a result, observations cannot provide enough evidence to support one of the various theoretical models of jet production over the many that exist. The jets have their most obvious observational effects in the radio waveband, where Very Long Baseline Interferometry can be used to study the synchrotron radiation they emit at resolutions of sub-parsec scales. However, they radiate in all wavebands from the radio through to the gamma-ray range via the synchrotron and the inverse-Compton scattering process, and so AGN jets are a second potential source of any observed continuum radiation.
Radiatively inefficient AGN
There exists a class of 'radiatively inefficient' solutions to the equations that govern accretion. The most widely known of these is the Advection Dominated Accretion Flow (ADAF), but other theories exist. In this type of accretion, which is important for accretion rates well below the Eddington limit, the accreting matter does not form a thin disc and consequently does not efficiently radiate away the energy that it acquired as it moved close to the black hole. Radiatively inefficient accretion has been used to explain the lack of strong AGN-type radiation from massive black holes at the centres of elliptical galaxies in clusters, where otherwise we might expect high accretion rates and correspondingly high luminosities. Radiatively inefficient AGN would be expected to lack many of the characteristic features of standard AGN with an accretion disc.
Observational characteristics
There is no single observational signature of an AGN. The list below covers some of the historically important features that have allowed systems to be identified as AGN.
- Nuclear optical continuum emission. This is visible whenever there is a direct view of the accretion disc. Jets can also contribute to this component of the AGN emission. The optical emission has a roughly power-law dependence on wavelength.
- Nuclear infra-red emission. This is visible whenever the accretion disc and its environment are obscured by gas and dust close to the nucleus and then re-emitted ('reprocessing'). As it is thermal emission, it can be distinguished from any jet or disc-related emission.
- Broad optical emission lines. These come from cold material close to the central black hole. The lines are broad because the emitting material is revolving around the black hole with high speeds causing a range of Doppler shifts of the emitted photons.
- Narrow optical emission lines. These come from more distant cold material, and so are narrower than the broad lines.
- Radio continuum emission. This is always due to a jet. It shows a spectrum characteristic of synchrotron radiation.
- X-ray continuum emission. This can arise both from a jet and from the hot corona of the accretion disc via a scattering process: in both cases it shows a power-law spectrum. In some radio-quiet AGN there is an excess of soft X-ray emission in addition to the power-law component. The origin of the soft X-rays is not clear at present.
- X-ray line emission. This is a result of illumination of cold heavy elements by the X-ray continuum that causes fluorescence of X-ray emission lines. The best-known of which is the iron feature around 6.4 keV. This line may be narrow or broad: relativistically broadened iron lines can be used to study the dynamics of the accretion disc very close to the nucleus and therefore the nature of the central black hole.
Types of active galaxy
It is convenient to divide AGN into two classes, conventionally called radio-quiet and radio-loud. In the radio-loud objects the emission contribution from the jet(s) and the lobes that they inflate dominates the luminosity of the AGN, at least at radio wavelengths but possibly at some or all others. Radio-quiet objects are simpler since jet and jet-related emission can be neglected.
AGN terminology is often confusing, since the distinctions between different types of AGN sometimes reflect historical differences in how the objects were discovered or initially classified, rather than real physical differences.
Radio-quiet AGN
- Low-ionization nuclear emission-line regions (LINERs). As the name suggests, these systems show only weak nuclear emission-line regions, and no other signatures of AGN emission. It is debatable whether all such systems are true AGN (powered by accretion on to a supermassive black hole). If they are, they constitute the lowest-luminosity class of radio-quiet AGN. Some may be radio-quiet analogues of the low-excitation radio galaxies (see below).
- Seyfert galaxies. Seyferts were the earliest distinct class of AGN to be identified. They show optical range nuclear continuum emission, narrow and occasionally broad emission lines, occasionally strong nuclear X-ray emission and sometimes a weak small-scale radio jet. Originally they were divided into two types known as Seyfert 1 and 2: Seyfert 1s show strong broad emission lines while Seyfert 2s do not, and Seyfert 1s are more likely to show strong low-energy X-ray emission. Various forms of elaboration on this scheme exist: for example, Seyfert 1s with relatively narrow broad lines are sometimes referred to as narrow-line Seyfert 1s. The host galaxies of Seyferts are usually spiral or irregular galaxies.
- Radio-quiet quasars/QSOs. These are essentially more luminous versions of Seyfert 1s: the distinction is arbitrary and is usually expressed in terms of a limiting optical magnitude. Quasars were originally 'quasi-stellar' in optical images as they had optical luminosities that were greater than that of their host galaxy. They always show strong optical continuum emission, X-ray continuum emission, and broad and narrow optical emission lines. Some astronomers use the term QSO (Quasi-Stellar Object) for this class of AGN, reserving 'quasar' for radio-loud objects, while others talk about radio-quiet and radio-loud quasars. The host galaxies of quasars can be spirals, irregulars or ellipticals. There is a correlation between the quasar's luminosity and the mass of its host galaxy, in that the most luminous quasars inhabit the most massive galaxies (ellipticals).
- 'Quasar 2s'. By analogy with Seyfert 2s, these are objects with quasar-like luminosities but without strong optical nuclear continuum emission or broad line emission. They are scarce in surveys, though a number of possible candidate quasar 2s have been identified.
Radio-loud AGN
See main article Radio galaxy for a discussion of the large-scale behaviour of the jets. Here, only the active nuclei are discussed.
- Radio-loud quasars behave exactly like radio-quiet quasars with the addition of emission from a jet. Thus they show strong optical continuum emission, broad and narrow emission lines, and strong X-ray emission, together with nuclear and often extended radio emission.
- “Blazars” (BL Lac objects and OVV quasars) classes are distinguished by rapidly variable, polarized optical, radio and X-ray emission. BL Lac objects show no optical emission lines, broad or narrow, so that their redshifts can only be determined from features in the spectra of their host galaxies. The emission-line features may be intrinsically absent or simply swamped by the additional variable component. In the latter case, emission lines may become visible when the variable component is at a low level. OVV quasars behave more like standard radio-loud quasars with the addition of a rapidly variable component. In both classes of source, the variable emission is believed to originate in a relativistic jet oriented close to the line of sight. Relativistic effects amplify both the luminosity of the jet and the amplitude of variability.
- Radio galaxies. These objects show nuclear and extended radio emission. Their other AGN properties are heterogeneous. They can broadly be divided into low-excitation and high-excitation classes. Low-excitation objects show no strong narrow or broad emission lines, and the emission lines they do have may be excited by a different mechanism. Their optical and X-ray nuclear emission is consistent with originating purely in a jet. They may be the best current candidates for AGN with radiatively inefficient accretion. By contrast, high-excitation objects (narrow-line radio galaxies) have emission-line spectra similar to those of Seyfert 2s. The small class of broad-line radio galaxies, which show relatively strong nuclear optical continuum emission probably includes some objects that are simply low-luminosity radio-loud quasars. The host galaxies of radio galaxies, whatever their emission-line type, are essentially always ellipticals.
These galaxies can be broadly summarised by the following table:
|Emission Lines||X-rays||Excess of||Strong
|OVV||yes||no||stronger than BL Lac||yes||yes||no||yes||yes||yes||yes|
Unification of AGN species
Unified models propose that different observational classes of AGN are really a single type of physical object observed under different conditions. The currently favoured unified models are 'orientation-based unified models' meaning that they propose that the apparent differences between different types of objects arise simply because of their different orientations to the observer.
Radio-quiet unification
At low luminosities, the objects to be unified are Seyfert galaxies. The unification models propose that in Seyfert 1s the observer has a direct view of the active nucleus. In Seyfert 2s the nucleus is observed through an obscuring structure which prevents a direct view of the optical continuum, broad-line region or (soft) X-ray emission. The key insight of orientation-dependent accretion models is that the two types of object can be the same if only certain angles to the line of sight are observed. The standard picture is of a torus of obscuring material surrounding the accretion disc. It must be large enough to obscure the broad-line region but not large enough to obscure the narrow-line region, which is seen in both classes of object. Seyfert 2s are seen through the torus. Outside the torus there is material that can scatter some of the nuclear emission into our line of sight, allowing us to see some optical and X-ray continuum and, in some cases, broad emission lines—which are strongly polarized, showing that they have been scattered and proving that some Seyfert 2s really do contain hidden Seyfert 1s. Infrared observations of the nuclei of Seyfert 2s also support this picture.
At higher luminosities, quasars take the place of Seyfert 1s, but, as already mentioned, the corresponding 'quasar 2s' are elusive at present. If they do not have the scattering component of Seyfert 2s they would be hard to detect except through their luminous narrow-line and hard X-ray emission.
Radio-loud unification
Historically, work on radio-loud unification has concentrated on high-luminosity radio-loud quasars. These can be unified with narrow-line radio galaxies in a manner directly analogous to the Seyfert 1/2 unification (but without the complication of much in the way of a reflection component: narrow-line radio galaxies show no nuclear optical continuum or reflected X-ray component, although they do occasionally show polarized broad-line emission). The large-scale radio structures of these objects provide compelling evidence that the orientation-based unified models really are true. X-ray evidence, where available, supports the unified picture: radio galaxies show evidence of obscuration from a torus, while quasars do not, although care must be taken since radio-loud objects also have a soft unabsorbed jet-related component, and high resolution is necessary to separate out thermal emission from the sources' large-scale hot-gas environment. At very small angles to the line of sight, relativistic beaming dominates, and we see a blazar of some variety.
However, the population of radio galaxies is completely dominated by low-luminosity, low-excitation objects. These do not show strong nuclear emission lines — broad or narrow — they have optical continua which appear to be entirely jet-related, and their X-ray emission is also consistent with coming purely from a jet, with no heavily absorbed nuclear component in general. These objects cannot be unified with quasars, even though they include some high-luminosity objects when looking at radio emission, since the torus can never hide the narrow-line region to the required extent, and since infrared studies show that they have no hidden nuclear component: in fact there is no evidence for a torus in these objects at all. Most likely, they form a separate class in which only jet-related emission is important. At small angles to the line of sight, they will appear as BL Lac objects.
Cosmological uses and evolution
For a long time, active galaxies held all the records for the highest-redshift objects known either in the optical or the radio spectrum, because of their high luminosity. They still have a role to play in studies of the early universe, but it is now recognised that an AGN gives a highly biased picture of the 'typical' high-redshift galaxy.
More interesting is the study of the evolution of the AGN population. Most luminous classes of AGN (radio-loud and radio-quiet) seem to have been much more numerous in the early universe. This suggests (1) that massive black holes formed early on and (2) that the conditions for the formation of luminous AGN were more common in the early universe, such as a much higher availability of cold gas near the centre of galaxies than at present. It also implies that many objects that were once luminous quasars are now much less luminous, or entirely quiescent. The evolution of the low-luminosity AGN population is much less well understood due to the difficulty of observing these objects at high redshifts.
See also
- Lynden-Bell, D. (1969). "Galactic Nuclei as Collapsed Old Quasars". Nature 223 (5207): 690–694. Bibcode:1969Natur.223..690L. doi:10.1038/223690a0.
- Kazanas, Demosthenes (2012). "Toward a Unified AGN Structure". Astronomical Review 7 (2).
- Marconi, A.; L. K. Hunt (2003). "The Relation between Black Hole Mass, Bulge Mass, and Near-Infrared Luminosity". The Astrophysical Journal 589 (1): L21–L24. arXiv:astro-ph/0304274. Bibcode:2003ApJ...589L..21M. doi:10.1086/375804.
- Narayan, R.; I. Yi (1994). "Advection-Dominated Accretion: A Self-Similar Solution". Astrophys. J 428: L13. arXiv:astro-ph/9403052. Bibcode:1994ApJ...428L..13N. doi:10.1086/187381.
- Fabian, A. C.; Rees; M. J. Rees (1995). "The accretion luminosity of a massive black hole in an elliptical galaxy". Monthly Notices of the Royal Astronomical Society 277 (2): L55–L58. arXiv:astro-ph/9509096. Bibcode:1995MNRAS.277L..55F.
- Vermeulen, R. C.; P. M. Ogle, H. D. Tran, I. W. A. Browne, M. H. Cohen, A. C. S. Readhead, G. B. Taylor, R. W. Goodrich (1995). "When Is BL Lac Not a BL Lac?". The Astrophysical Journal Letters 452 (1): 5–8. Bibcode:1995ApJ...452L...5V. doi:10.1086/309716.
- HINE, RG; Longair; MS LONGAIR (1979). "Optical spectra of 3 CR radio galaxies". Royal Astronomical Society, Monthly Notices 188: 111–130. Bibcode:1979MNRAS.188..111H.
- Laing, R. A.; C. R. Jenkins, J. V. Wall, S. W. Unger (1994). "Spectrophotometry of a Complete Sample of 3CR Radio Sources: Implications for Unified Models". The First Stromlo Symposium: the Physics of Active Galaxies. ASP Conference Series, 54.
- Baum, S. A.; E. L. Zirbel, C. P. O'Dea (1995). "Toward Understanding the Fanaroff-Riley Dichotomy in Radio Source Morphology and Power". The Astrophysical Journal 451: 88. Bibcode:1995ApJ...451...88B. doi:10.1086/176202.
- Chiaberge, M.; A. Capetti, A. Celotti (2002). "Understanding the nature of FRII optical nuclei: a new diagnostic plane for radio galaxies". Journal reference: Astron. Astrophys 394 (3): 791–800. arXiv:astro-ph/0207654. Bibcode:2002A&A...394..791C. doi:10.1051/0004-6361:20021204.
- Hardcastle, M. J.; D. A. Evans, J. H. Croston (2006). "The X-ray nuclei of intermediate-redshift radio sources". Monthly Notices of the Royal Astronomical Society 370 (4): 1893–1904. arXiv:astro-ph/0603090. Bibcode:2006MNRAS.370.1893H. doi:10.1111/j.1365-2966.2006.10615.x.
- Grandi, S. A.; D. E. Osterbrock (1978). "Optical spectra of radio galaxies". Astrophysical Journal 220 (Part 1): 783. Bibcode:1978ApJ...220..783G. doi:10.1086/155966.
- Antonucci, R. (1993). "Unified Models for Active Galactic Nuclei and Quasars". Annual Reviews in Astronomy and Astrophysics 31 (1): 473–521. Bibcode:1993ARA&A..31..473A. doi:10.1146/annurev.aa.31.090193.002353.
- Urry, P.; Paolo Padovani (1995). "Unified schemes for radioloud AGN". Publications of the Astronomical Society of the Pacific 107: 803–845. arXiv:astro-ph/9506063. Bibcode:1995PASP..107..803U. doi:10.1086/133630.
- Laing, R. A. (1988). "The sidedness of jets and depolarization in powerful extragalactic radio sources". Nature 331 (6152): 149–151. Bibcode:1988Natur.331..149L. doi:10.1038/331149a0.
- Garrington, S. T.; J. P. Leahy, R. G. Conway, RA LAING (1988). "A systematic asymmetry in the polarization properties of double radio sources with one jet". Nature 331 (6152): 147–149. Bibcode:1988Natur.331..147G. doi:10.1038/331147a0.
- Barthel, P. D. (1989). "Is every quasar beamed?". Astrophysical Journal 336: 606–611. Bibcode:1989ApJ...336..606B. doi:10.1086/167038.
- Belsole, E.; D. M. Worrall, M. J. Hardcastle (2006). "High-redshift Faranoff-Riley type II radio galaxies: X-ray properties of the cores". Monthly Notices of the Royal Astronomical Society 366 (1): 339–352. arXiv:astro-ph/0511606. Bibcode:2006MNRAS.366..339B. doi:10.1111/j.1365-2966.2005.09882.x.
- Ogle, P.; D. Whysong, R. Antonucci (2006). "Spitzer Reveals Hidden Quasar Nuclei in Some Powerful FR II Radio Galaxies". The Astrophysical Journal 647 (1): 161–171. arXiv:astro-ph/0601485. Bibcode:2006ApJ...647..161O. doi:10.1086/505337.
- Browne, I. W. A. (1983). "Is it possible to turn an elliptical radio galaxy into a BL Lac object?". Monthly Notices of the Royal Astronomical Society 204: 23P–27P. Bibcode:1983MNRAS.204P..23B.
- Media related to Active galactic nuclei at Wikimedia Commons |
A long but worthy read. Sadly.
An Overview of the Experiences of LGBT Youth in the Juvenile Justice System
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Gay, transgender, and gender nonconforming youth are significantly over-represented in the juvenile justice system—approximately 300,000 gay and transgender youth are arrested and/or detained each year, of which more than 60 percent are black or Latino. Though gay and transgender youth represent just 5 percent to 7 percent of the nation’s overall youth population, they compose 13 percent to 15 percent of those currently in the juvenile justice system.
These high rates of involvement in the juvenile justice system are a result of gay and transgender youth abandonment by their families and communities, and victimization in their schools—sad realities that place this group of young people at a heightened risk of entering the school-to-prison pipeline.
Despite the disproportionately high rates of gay and transgender youth entering the juvenile justice system, our nation’s schools, law enforcement officers, district attorneys, judges, and juvenile defenders are not equipped to manage the unique experiences and challenges that these young people face. As a consequence, the system often does more harm by unfairly criminalizing these youth—imposing harsh school sanctions, labeling them as sex offenders, or detaining them for minor offenses—in addition to subjecting them to discriminatory and harmful treatment that deprives them of their basic civil rights.
Angela Irvine of the National Council on Crime and Delinquency in conjunction with the Equity Project, which works to ensure gay and transgender youth in the juvenile justice system are treated with fairness and respect, have both generated groundbreaking research on the experiences of these youth in the system over the past few years. This issue brief offers a high-level summary of some of their findings, as well as others, to explain the disproportionate pipelining of gay and transgender youth into the juvenile justice system, the bias and discrimination they face once within the system, and the steps that the federal government and state and local juvenile courts can take to ensure that gay and transgender youth are treated with dignity and respect.
Why gay and transgender youth end up in the juvenile justice system
Family rejection, homelessness, and failed safety nets
Research shows that gay and transgender youth entering into the juvenile justice system are twice as likely to have experienced family conflict, child abuse, and homelessness as other youth. This trend is partly due to the fact that youth today “come out” at younger ages, often to families that may not accept gay and transgender people. Since these youth still depend on their families to meet their material needs, family rejection can leave them emotionally and physically vulnerable, particularly if they find themselves cast onto the streets with nowhere to turn for support.
Many gay and transgender youth leave their homes of their own accord to escape the conflict and emotional or physical abuse that can ensue—26 percent report leaving their homes at some point—but more often, they are pushed out and into the juvenile justice system by their own families.
Interfamily conflicts stemming from parents’ refusal to accept a child’s sexual orientation or gender identity often result in the first contact these young people have with the justice system. According to the Equity Project, prosecutors frequently file charges against these youth for being “incorrigible” or beyond the control of their parents or guardians, based largely on the parent’s objections to their sexual orientation. This practice unfairly criminalizes gay and transgender youth because of their identity rather than because of their behavior.
Further, family discord that casts these youth from their homes can send them cascading through social safety nets not adequately equipped to support them. Programs designed to keep children and youth off the streets, such as foster care, health centers, and other youth-serving institutions, are often ill-prepared or unsafe for gay and transgender youth due to institutional prejudice, lack of provider and foster-parent training, and discrimination against gay and transgender youth by adults and peers.As a result, many youth run away from these placements, actions that could also land them in the custody of the juvenile justice system.
Gay and transgender youth who flee hostility and abuse at home and in temporary placements are most likely to end up homeless, which is the greatest predictor of involvement with the juvenile justice system. Gay and transgender youth represent up to 40 percent of the homeless youth population even though they only compose 5 percent to 7 percent of the youth population overall, and 39 percent of homeless gay and transgender youth report being involved in the juvenile justice system at some level.
Out of despair and a need for survival, homeless gay and transgender youth are more likely to resort to criminal behaviors, such as drug sales, theft, or “survival sex,” which put them at risk of arrest and detainment. These youth are also at an increased risk of detainment for committing crimes related to homelessness, such as violating youth curfew laws and sleeping in public spaces.
Family rejection, which sets off a tragic chain of events for many gay and transgender youth, is at the core of these issues. Caitlin Ryan of the Family Acceptance Project at San Francisco State University, whose research has brought to light the negative impacts that family rejection can have on gay and transgender youth, emphasizes the need to provide opportunities to help support and strengthen families in order to promote nurturing environments for gay and transgender children. Early intervention can help families and caregivers reduce the risk of these youth entering the juvenile justice system. It is important that law enforcement officials, district attorneys, judges, and juvenile defenders seek ways to keep gay and transgender youth and their families together, rather than pushing for incarceration.
Biased school discipline policies
Unfortunately, schools do not always provide a reprieve for youth experiencing family rejection. According to the Gay Lesbian and Straight Education Network’s School Climate Survey, 84 percent of gay and transgender students report being verbally harassed, 40 percent physically harassed, and 19 percent physically assaulted.
What’s more, gay and transgender students report astonishingly low levels of confidence in their school administrators and often do not report incidents because they expect the situation will not improve or fear it might even become worse. This is not surprising considering that one-third of bullied gay and transgender students who reported bullying to school officials said the administrators did nothing to address the issue.
In fact, school officials in many ways exacerbate these problems and place further stress and burden on gay and transgender youth by disproportionately doling out harsh school sanctions against them for minor disciplinary infractions. The school and juvenile justice systems have become inextricably linked in recent years with schools relying heavily on law enforcement to manage what in the past were school discipline issues. The consequence of this conflated discipline system is that it unduly criminalizes youth of color and gay and transgender youth.
School discipline policies across the United States are under heightened scrutiny because of the disparate impact they have on youth of color, particularly black boys. Data released this spring from the U.S. Department of Education’s Office of Civil Rights show that harsh school sanctions—such as zero-tolerance policies, which lead to suspensions and expulsions of students for even the most minor offenses—perpetuate a school-to-prison pipeline that disproportionately criminalizes youth of color.
Hidden among these school discipline data are thousands of gay and transgender youth who bear a double burden of disparate impact. A groundbreaking study published in 2010 in the medical journal Pediatrics revealed that gay and transgender youth, particularly gender nonconforming girls, are up to three times more likely to experience harsh disciplinary treatment by school administrators than their heterosexual counterparts.
As with the racial disparities in school suspensions and expulsions, these higher rates of punishment do not correlate to higher rates of misbehavior among gay and transgender youth. What the research suggests is that gay and transgender youth actually face harsher sanctions by school administrators even when committing similar offenses.
Surely bias and discrimination among teachers, staff, and administrators contributes to the unfair treatment of gay and transgender youth in schools. Adults in schools often draw assumptions of guilt based on a student’s physical characteristics, demeanor, dress, or mannerisms, deeming those deviating from an accepted gender norm to be agitators. Such assumptions are not only misguided, but biased against gay and transgender students who do not fall within rigid stereotypes of expression.
Moreover, studies reveal that gay and transgender youth are often the victims, rather than the aggressors in school conflicts, which stem from bullying and harassment. Consider, for example, a gender nonconforming girl exhibiting masculine traits, who is disciplined for fighting but may be defending herself from peers’ taunts. Yet more often than not, school administrators will consider her the aggressor based solely on her physical demeanor and will suspend or expel her despite the defensive nature of her actions.
For many students, suspension and expulsion are the first steps toward time behind bars. This is equally true for gay and transgender youth. Black boys and gender nonconforming girls similarly experience disproportionately harsh punishments and juvenile justice system referrals in schools, but the latter are rendered all but invisible because sexual orientation and gender identity are not included in the federal school discipline data cited earlier in this report. A first step in addressing the unfair punishment of gay and transgender youth in schools is to expand the research and collection of school discipline data to include gay and transgender youth, which will help policymakers and practitioners alike better understand the problem and formulate more supportive school discipline policies.
Unfair criminalization by the system
Once in the juvenile justice system, gay and transgender youth are too often denied basic civil rights, wrongly categorized as sexually deviant simply because of their sexual orientation, gender identity, or gender nonconformity, and even labeled as sex offenders. They are also subjected to the biases and discrimination of law enforcement agents, judges, and other justice system officials that leave them vulnerable to abuse and neglect.
Classification as sex offenders
Gay and transgender youth who end up in the justice system are at-risk of being labeled as sex offenders, regardless of whether they have actually committed a sexual crime. Gay and transgender youth “are more likely to be prosecuted for age-appropriate consensual sexual activity” than their heterosexual counterparts—a lopsided application of the law, which has devastating consequences for gay and transgender youth who would be required to register as a sex offenders in 29 states if convicted. The stigma of being a registered sex offender could haunt them for the rest of their lives, negatively impacting their future employment and life opportunities and causing significant psychological distress.
Many gay and transgender youth charged with nonsexual offenses are also unfairly treated as sex offenders and ordered by the court to undergo sex offender treatment programs or sex offense risk assessments simply because of their sexual orientation or gender identity. This misguided categorization by the courts has led gay and transgender youth, innocent of violent crimes or sex offenses, to be placed in restrictive punitive settings for high-risk youth and to be given longer stays in out-of-home placements.
These restrictive settings not only hinder rehabilitation efforts, they perpetuate the stigma that being gay or transgender is wrong. Additionally, extended stays in out-of-home placements prevent gay and transgender youth from reconnecting with their families, a critical step proven to stabilize their lives and reduce their risk of returning to the system. These unfair practices make gay and transgender youth susceptible to discrimination and harmful treatment while in the system.
Detention as a default
In most incidences juveniles who have been arrested or detained will only be released from custody under the supervision of a parent or guardian. Without someone to claim them, youth can be left to languish in detention centers with youth convicted of crimes, even if they have not been. Gay and transgender youth are most at-risk of detainment by default by the juvenile justice system as they are more likely to be estranged from their families and lack parental support, which leaves them to fend for themselves. As a consequence, these youth are subjected to criminal incarceration while they await foster or group home placements.
Discriminatory and harmful treatment
Segregation and isolation of gay and transgender youth
From the moment gay and transgender youth enter a detention facility they are at risk of being inappropriately classified and housed. Transgender youth, for example, are often placed according to their birth sex rather than by their gender identity in an effort to force transgender youth to conform to societal norms. Doing so can be psychologically devastating and leave them vulnerable to physical and sexual abuse. Additionally, youth facility staff often view them as threatening or sexually predatory, harmful stereotypes that taint placement decisions and influence the treatment of transgender youth.
Some facilities will automatically segregate gay and transgender youth or place them in solitary confinement for their “own safety,” but this isolation perpetuates the stigmatization of gay and transgender youth, casts them as sexually deviant, and signals that they might be of threat to other youth.
According to the American Psychiatric Association, isolation “is a form of punishment and is likely to produce lasting psychiatric symptoms.” Unwarranted segregation deprives gay and transgender youth of educational, recreational, and programming opportunities that they are otherwise entitled to receive, punishing them unfairly and at a particularly vulnerable time in their adolescent development.
Physical, sexual, and emotional abuse
A 2007 study funded by the California Department of Corrections and Rehabilitation found an astounding 67 percent of gay or transgender men have been sexually assaulted by another inmate—a rate 15 times higher than the overall inmate population. Another study found that sexual assaults that occur are not just isolated events, but that 30 percent of all inmates have endured six or more sexual assaults.
Gay and transgender youth are particularly at risk for physical, sexual, and emotional abuse while in detention, by both staff and other youth. Eighty percent of those surveyed by the Equity Project believed a lack of safety in dentition was a serious problem. Some reports suggest that staff have turned a blind eye to incidents of rape and abuse against gay and transgender youth, confusing gay and transgender identity as an invitation for sex. Gay and transgender youth are not only subjected to abuse by their peers but by staff as well, particularly in the facilities that lack training and policies that promote inclusiveness and rely on biases rather than on best practices in treatment and placement decisions. This type of environment allows physical, sexual, and emotional abuse toward gay and transgender youth to happen without so much as a second thought and leaves them with nowhere to turn for help.
Unsafe reparative or conversion therapy
Gay and transgender youth have been subjected to reparative or conversion therapy to change their sexual orientation by both social workers and the courts, even though so-called reparative or conversion therapy has been condemned by every major health organization, including the American Medical Association, American Psychological Association, and the American Academy of Child and Adolescent Psychiatry.
Sadly, the juvenile justice system is rife with examples of misguided interventions. One judge hospitalized a gay youth to stop his same-sex attraction, while another judge with the parent’s approval, had a young lesbian who was caught in a sexual act with another girl placed in a private hospital to be “treated and diagnosed for this behavior.” These examples may be the extreme, but instances such as a 15-year-old boy being given a women’s lingerie catalogue with the purpose of teaching him “appropriate” sexual desires and a male-to-female transgender youth, who was detained in a boy’s facility, being placed on “treatment plan” to “help with gender confusion and appropriate gender identity,” are more common examples of unsafe reparative therapy.
The inclination to change a youth’s sexual orientation or gender identity or force him or her to conform to “social norms” hinders general mental health and causes severe psychological distress. This type of “counseling and other services are virtually worthless [for gay and transgender youth] because they either ignore or criminalize the youth’s sexuality.”
Conclusion and recommendations
Gay and transgender youth are pipelined into the juvenile justice system at disproportionate rates, often stripped of their basic dignity and civil rights, and treated in a harmful and discriminatory manner once in the system. The current policies and practices of schools and the juvenile justice system overlook gay and transgender youth and perpetuate stigma and bias that can lead to their unwarranted criminalization and unfair treatment.
Some of the issues discussed in this report stem from the lack of cultural competency on the part of school officials, law enforcement officers, district attorneys, judges, and juvenile defenders. The individuals who interface directly with these youth must be better equipped to provide respectful, culturally appropriate interventions in order to reduce the number of gay and transgender youth unfairly and unnecessarily pipelined into the juvenile justice system and to improve conditions for them once in the system. A first step toward improving the system will be to institute training standards across the board for all agents of the court.
Moreover, institutional bias is at the heart of the mistreatment of gay and transgender youth by schools and the juvenile justice system, and we recommend broad policy suggestions to address them. The following recommendations are adopted in part from the National Juvenile Justice and Delinquency Prevention Council and the Equity Project. These recommendations are by no means the only recommendations for improvement, but instead offer a start to address the serious issue of the criminalization of gay and transgender youth.
- Promotion of family center interventions
Family rejection drives many gay and transgender youth from their homes and perpetuates negative coping behaviors and unlawful activities, heightening risk of entering the juvenile justice system. Gay and transgender youth who feel a sense of family acceptance report better physical, mental, and educational outcomes all around. Yet families have largely been left out of the equation on reform. Interventions that reconnect youth with their families will reduce their susceptibility to involvement with the juvenile justice system. Moreover, schools and juvenile detention systems should engage in more supportive behaviors that reduce risk and promote the positive development of gay and transgender youth in custodial care.
- Gay and transgender inclusive training for all juvenile justice professionals
“Juvenile justice professionals must receive training and resources regarding the unique societal, familial, and developmental challenges confronting [gay and transgender] youth … [T]rainings must be designed to address the specific professional responsibilities of the audience (i.e., judges, defense attorneys, prosecutors, probation officers, and detention staff).” Furthermore, these professionals must ensure that they and others treat gay and transgender youth with dignity, respect, and fairness and avoid ridiculing or attempting to change their sexual orientation or gender identity.
- Development of gay and transgender inclusive policies, procedures, and programs
All agencies involved with juvenile justice and all officers of the juvenile court should develop, adopt, and enforce polices that prohibit discrimination and mistreatment of any youth on the basis of their actual or perceived sexual orientation or gender identity at all stages of the process. Juvenile justice professionals should develop appropriate responses to the behavior of each gay and transgender youth that are tailored to address specific needs and to promote individual well-being by allowing gay and transgender youth to express themselves freely. This would include giving them the choice of name, clothing, hairstyle, or any other means by which they feel comfortable expressing themselves.
- Gay and transgender inclusive data collection by the Office of Juvenile Justice and Delinquency Prevention
Many states and localities are unable to achieve meaningful changes to their juvenile justice systems because the lack adequate data. Therefore, the U.S. Department of Justice’s Office of Juvenile Justice and Delinquency Prevention should prioritize data collection that is disaggregated not only by race, ethnicity, and gender, but also by sexual orientation and gender identity. Through this type of robust data collection, communities will be better able to develop services that are culturally and linguistically appropriate for youth and their families, especially for gay and transgender youth.
- Dismantle the school-to-prison pipeline for all youth and for gay and transgender youth in particular
As previously noted, gay and transgender youth are disproportionately pipelined into the juvenile justice system by their schools. Yet little data are available to quantify the problem and address it effectively. We can address this by including sexual orientation and gender identity questions in the Department of Education’s school discipline data collection efforts. In addition, it is important that the Supportive School Discipline Initiative, a joint effort between the Department of Education and the Department of Justice aiming to dismantle the school-to-prison pipeline and currently focusing on race, provide guidance to school systems that bridges the gap for all youth, including those who are gay or transgender.
- Pass the Safe Schools Improvement Act and the Student Nondiscrimination Act
The federal Student Nondiscrimination Act would prohibit discrimination on the basis of gender identity and sexual orientation against any student in a public school that receives federal funding and would allow individuals to take legal action and be awarded compensatory damages and reimbursement of court costs if judgment is found in their favor under the bill’s provisions. The Safe Schools Improvement Act would require kindergarten-through-12th-grade public schools that receive federal funding to implement policies prohibiting harassment and bullying based on gender identity and sexual orientation. The bill would also require states to report harassment and bullying data to the U.S. Department of Education. Passage of these two pieces of legislation by Congress would help ensure that gay and transgender students do not end up in the juvenile justice system for protecting themselves against discrimination they face on a daily basis at school.
- Gay and transgender cultural competence in Safe Schools/Healthy Students
The Department of Education’s Safe Schools/Healthy Students program is widely recognized as a model of “effective collaboration across public education, local mental health, and juvenile justice.” Therefore the evaluations of the Safe Schools/Healthy Students program should be inclusive of gay and transgender students needs, including cultural competency training for those who work in the juvenile justice system. Increased training focused on the needs of gay and transgender youth will help decrease arrest rates and referrals of these youth to the juvenile justice system. Currently, law enforcement is 50 percent more likely to stop gay youth than other youth, according to a recent study in the medical journal Pediatrics. Additionally, “girls who identified themselves as lesbians or bisexual reported twice as many arrests and convictions as other girls who had engaged in similar behavior.” Gay and transgender youth in general experience high rates of school violence, which not only interferes with their ability to learn but also can affect their involvement with the juvenile justice system.
- Amend the Sex Offender Registration and Notification Act
Congress should amend the Sex Offender Registration and Notification Act, or SORNA, to exclude youth who are convicted of certain sex-based offenses from mandatory sex offender registration. The way the U.S. Department of Justice’s SORNA program is currently set up can disrupt families because both the youth who has to register is stigmatized and the entire family. More importantly, SORNA can have lasting effects on a youth’s identification and treatment. A parent, for example, may delay getting his or her child needed treatment and even go as far as hiding the child’s problem after discovering that the child may have to register as a sex offender for life. Thus, amending this law can help to alleviate the criminalization that is part and parcel of a gay or transgender youth being labeled as a sex offender.
- Reauthorization of the Juvenile Justice and Delinquency Prevention Act
The Juvenile Justice and Delinquency Prevention Act has not been reauthorized since 2002 and is in need of substantive changes to ensure that youth are put on a better path and that our communities are kept safe. In particular, the act should be reauthorized to include three key items that will help gay and transgender youth in the system. The first is to strengthen the Disproportionate Minority Contact core protections “by requiring states to take concrete steps to reduce racial and ethnic disparities in the juvenile justice system,” rather than just issuing vague requirements with no clear guidance on how to reduce these disparities.
Second, the partnership between states and the Office of Juvenile Justice and Delinquency Prevention should expand training, technical assistance, and research and evaluation to include of gay and transgender needs and issues.
Third, incentives for the juvenile justice system to ensure that all policies, practices, and programs recognize the unique needs of gay and transgender youth, including accountability measures, expertise on the Juvenile Justice and Delinquency Prevention Act advisory groups, and increased research and information dissemination should be created.
- Passage of federal legislation prohibiting gay and transgender discrimination in the juvenile justice system
As this brief has shown, the juvenile justice system holds a disproportionally large number of gay and transgender youth, who experience high rates of discrimination and violence. Passage of federal protections that would ensure equality and the end of discrimination against gay and transgender youth is universal and should not be legislated on a state-by-state basis. Therefore, “Congress should pass federal protections against discrimination in all settings based on actual or perceived sexual orientation or gender identity and create incentives for States to appropriately and effectively respond to [gay and transgender] youth involved in the justice system.”
Jerome Hunt is a Research Associate for LGBT Progress at American Progress and Aisha C. Moodie-Mills is the Advisor for LGBT Policy & Racial Justice at American Progress.
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|Type||Main battle tank|
|Place of origin||USSR|
|In service||1963 – present|
|Used by||Soviet Union, Belarus, Russia, Ukraine, Uzbekistan|
|Designer||Morozov Design Bureau|
|Designed||1951 – 62|
|Produced||1963 – 87|
|Weight||38 tonnes (42 short tons; 37 long tons)|
|Length||9.225 m (30 ft 3.2 in) (gun forward)|
|Width||3.415 m (11 ft 2.4 in)|
|Height||2.172 m (7 ft 1.5 in)|
|Armour||20–450 mm (0.79–18 in) of Glass-reinforced plastic sandwiched between layers of steel.
ERA plates optional
|D-81T 125 mm smoothbore gun|
|7.62 mm PKMT coax machine gun, 12.7 mm NSVT antiaircraft machine gun|
|Engine||5DTF 5-cyl. diesel
700 hp (522 kW)
|Power/weight||18.4 hp/tonne (13.7 kW/ton)|
|500 km (310 mi), 700 km (430 mi) with external tanks|
|Speed||60.5 km/h (37.6 mph) (road)|
The T-64 is a Soviet main battle tank, introduced in the early 1960s. It was a more advanced counterpart to the T-62: the T-64 served tank divisions, while the T-62 supported infantry in motor rifle divisions. Although the T-62 and the famed T-72 would see much wider use and generally more development, it was the T-64 that formed the basis of more modern Soviet tank designs like the T-80.
The T-64 was conceived in Kharkiv, Ukraine as the next-generation main battle tank by Alexander A. Morozov, the designer of the T-54 (which in the meantime would be incrementally improved by Leonid N. Kartsev's Nizhny Tagil bureau, in models T-54A, T-54B, T-55, and T-55A).
A revolutionary feature of the T-64 is the incorporation of an automatic loader for its 125-mm gun, allowing a crewmember's position to be omitted, and helping to keep the size and weight of the tank down. Tank troopers would joke that the designers had finally caught up with their unofficial hymn, "Three Tankers"—the song had been written to commemorate the crewmen fighting in the Battle of Khalkhin Gol, in 3-man BT-5 tanks in 1939.
The T-64 also pioneered other Soviet tank technology: the T-64A model of 1967 introduced the 125-mm smoothbore gun, and the T-64B of 1976 would be able to fire a guided antitank missile through its gun barrel.
The T-64 design was further developed as the gas turbine-powered T-80 main battle tank. The turret of the T-64B would be used in the improved T-80U and T-80UD, and an advanced version of its diesel engine would power T-80UD and T-84 tanks built in Ukraine.
The T-64 would be used only by the Soviet Army and never exported, unlike the T-54/55. It was superior to these tanks in most qualitative terms, until the introduction of the T-72B model in 1985. The tank equipped elite and regular formations in Eastern Europe and elsewhere, the T-64A model being first deployed with East Germany's Group of Soviet Forces in Germany (GSFG) in 1976, and some time later in Hungary's Southern Group of Forces (SFG). By 1981 the improved T-64B began to be deployed in East Germany and later in Hungary. While it was believed that the T-64 was "only" reserved for elite units, it was also used by much lower "non-ready formations", for example, the Odessa Military District's 14th Army.
With the break-up of the Soviet Union in 1991, T-64 tanks remained in the arsenals of constituent republics. Currently, slightly fewer than 2,000 of the old Soviet inventory of T-64 tanks are in service with the military of Ukraine and about 4,000 remain in service with the Russian Ground Forces.
Development history
The initial requirement
Recognizing that the T-55/T-62 lineage had finally exhausted its potential for improvement, the USSR embarked upon the development of an entirely new tank design that could defeat new Western tanks like the British Chieftain and resist new Western anti-tank weapons.
Project 430
Studies for the design of a new battle tank started as early as 1951. The KB-60M team was formed at the Kharkiv design bureau of the Kharkiv transport machine-building factory No. 75 named for Malyshev (Russian: конструкторское бюро Харьковского завода транспортного машиностроения №75 им. Малышева) by engineers coming back from Nizhniy Tagil, with A.A. Morozov at its head. A project named obyekt 430 gave birth to three prototypes which were tested in Kubinka in 1958. Those vehicles showed characteristics which were going to radically change the design of battle tanks on this side of the Iron Curtain. For the first time, an extremely compact opposed-piston engine was used : the 4TD, designed by the plant's engine design team. The transmission system comprised two lateral gears on each side of the engine. Those two innovations yielded a very short engine compartment with the opening located beneath the turret. The engine compartment volume was almost half that of the T-54. The cooling system was extracting and a new lightweight suspension was fitted, featuring hollow metallic wheels of a small diameter and caterpillar tracks with rubber joints.
The tank would keep a D-10TS 100 mm gun and frontal armour of 120 mm. As it did not present a clear superiority in terms of combat characteristics when compared to the T-55 which was entering active service, Morozov decided that production was not yet ready given the project's drawbacks. However, studies conducted on the obyekt 430U, featuring a 122 mm gun and 160 mm of armour, demonstrated that the tank had the potential to fit the firepower and armour of a heavy tank on a medium tank chassis. A new project was consequently started, obyekt 432.
Project 432
The gun fitted on this new tank was a powerful 115 mm D-68 (2A21). This was a potentially risky decision to replace the human loader by an electro-hydraulic automatic system, since the technology was new to Russian designers. The crew was reduced to three, which allowed an important reduction in internal volume, and consequently in weight, from 36 tonnes (obyekt 430) to 30.5 tonnes. The height dropped by 76 mm.
However, the arrival of the British 105 mm L7 gun and the US M-68 variant of it, fitted respectively to the Centurion and M60 tanks, forced the team to undertake another audacious première, with the adoption of a composite armour. The recently created process was called "K combination" by Western armies: this protection consisted of an aluminium alloy layer between two high strength steel layers. As a consequence, the weight of the prototype rose eventually to 34 tonnes. But as the engine was now a 700 hp (515 kW) 5TDF (also locally designed), its mobility remained excellent, far superior to the active T-62. The obyekt 432 was ready in September 1962 and the production started in October 1963 in Kharkiv plant. On December 30, 1966, it entered its service as the T-64.
Even as the first T-64s were rolling off the assembly lines, the design team was working on a new version which would allow it to keep firepower superiority, named obyekt 434. The brand new and very powerful 125 mm D-81T gun, from the Perm weapons factory, was fitted to the tank. This gun was merely a scaled-up version of the 115 mm smoothbore cannon from the T-62. The larger size of the 125 mm ammunition meant less could be carried inside the T-64, and with a fourth crewman loader taking up space as well, the tank would only have a 25-round capacity. This was unacceptably low for the Soviet designers, but strict dimensional parameters forbade them from enlarging the tank to increase interior space. The solution was to replace the human loader with a mechanical autoloader, cutting the crew to three and marking the first use of autoloaders in a Soviet MBT.(Perrett 1987:42) The 6ETs10 autoloader has 28 rounds and can fire 8 shots per minute; the stabiliser, a 2E23, was coupled to the new TPD-2-1 (1G15-1) sight. Night driving was also adapted with the new TPN-1-43A periscope which would benefit from the illumination of a powerful infrared L2G projector, fitted on the left side of the gun. The shielding was improved, with fibreglass replacing the aluminium alloy in the armour, and small spring-mounted plates fitted along the mudguards (known as the Gill skirt), to cover the top of the suspension and the side tanks. They were however extremely fragile and were often removed. Some small storage spaces were created along the turret, with a compartment on the right and three boxes on the front left. Schnorkels were mounted on the rear of the turret. A NBC protection system was fitted and the hatches were widened.
Prototypes were tested in 1966 and 1967 and, as production began after the six hundredth T-64, it entered service in the Soviet Army under the T-64A designation. Chief engineer Alexander Morozov was awarded the Lenin Prize for this model's success.
Designed for elite troops, the T-64A was constantly updated as available equipment was improved. After only three years in service, a first modernisation occurred, regarding :
- fire control, by replacing the sights with the TPD-2-49 day sight with an optical coincidence rangefinder and a TPN-1-49-23 night sight, and stabilisation by mounting a 2E26 system.
- the radio by mounting a R-123M
- night vision with a TBN-4PA for the driver and a TNP-165A for the tank leader. His battlepost was transformed by mounting a small stabilised turret with an anti-aircraft NSVT 12.7 mm x108 machine gun, electrically guided through an optical PZU-5 sight, and fed with 300 rounds. It could be used from within the tank so that the tank leader could avoid being exposed (as on previous tanks). The possibility of mounting a KMT-6 anti-mine system was also added.
A derived version appeared at the same time, designed for the commanding officer, and named T-64AK. It comprised a R-130M radio with a 10 m telescopic antenna which could be used only in a static position as it required shrouds, an artillery aiming circle PAB-2AM and TNA-3 navigation station, all of those could be supplied by an auxiliary gasoline-fired generator.
In 1976, the weapons system was improved by mounting a D-81TM (2A46-1), stabilised by a 2E28M2, supplied by an automatic 6ETs10M. The night sight is replaced by a TNPA-65 and the engine can accept different fuels, including diesel fuel, kerosene or gasoline. The production, first carried on the B variant, stopped in 1980.
But the majority of T-64As were still modernised after 1981, by mounting a six smoke grenade-launcher 81 mm 902A on each side of the gun, and by replacing the gill plates by a rubber skirt for a longer life. Some of them seem to have been fitted after 1985 with reactive bricks (as the T-64BV), or even with laser TPD-K1 telemeters instead of the optical TPD-2-49 optical coincidence rangefinder(1981). Almost all T-64's were modernised into T-64R, between 1977 and 1981, by reorganising external storage and snorkels, similar to the T-64A.
The design team was carrying on its work on new versions. Problems with the setup of the 5TDF engine occurred as the local production capacity was proven to be insufficient against a production done in three factories (Malyshev in Kharkiv, Kirov in Leningrad and Uralvagonzavod).
From 1961, an alternative to the obyekt 432 was studied, with 12 V-cylinder V-45 engine : the obyekt 436. Three prototypes were tested in 1966 in the Chelyabinsk factory. The order to develop a model derived from the 434 with the same engine gave the obyekt 438, later renamed as obyekt 439. Four tanks of this type were built and tested in 1969, which showed the same mobility as the production version, but mass production was not started. They served however as a basis for the design of the T-72 engine compartment.
In the beginning of the 1970s, the design team was trying to improve the tank further. The T-64A-2M study in 1973, with its more powerful engine and its reinforced turret, served as a basis for two projects :
- Obyekt 476 with a 6TD 1000 hp (735 kW) engine which served as a model for the T-80 combat compartment.
- Obyekt 447 which featured a new fire control with a laser telemeter, and which was able to fire missiles through the gun.
For the latter, the order was given to start its production under the name T-64B, as well as a derived version (which shared 95% of its components), the obyekt 437, without the missile guidance system for cost reasons. The latter was almost twice as much produced under the designation T-64B1. On September 3, 1976, the T-64B and the T-64B1 were declared good for the service, featuring the improved D-81Tm gun (2A46-2) with a 2E26M stabiliser, a 6ETs40 loader and a 1A33 fire control, including:
- a 1V517 ballistic calculator
- a 1G21 sight with laser telemetry
- a 1B11 cross-wind sensor.
Its ford capacity reaches 1.8 m without equipment. The T-64B had the ability to fire the new 9M112 "Kobra" radio-guided missile (NATO code "AT-8 Songster"). The vehicle then carries 8 missiles and 28 shells. The missile control system is mounted in front of the tank leader small turret and has many changes. The T-64B1 carries only 37 shells and has 2,000 7.62 mm rounds, against 1,250 for the T-64B.
They were modernised in 1981 by the replacement of the gun by a 2A46M1, the stabiliser by a 2E42, and the mounting of a 902A "Tucha-1" smoke grenade launcher in two groups of four, on each side of the gun. Two command versions are realised, very similar to the T-64AK: the T-64BK and the T-64B1K.
The decision, in October 1979, to start the production of the 6TD engine, and its great similarity with the 5TDF engine, allowed after some study to fit it in versions B and B1, but also A and AK, yielding the new models T-64AM, T-64AKM, T-64BM and T-64BAM, entering service in 1983.
The production ended in 1987 for all versions. The total production has reached almost 13,000.
Modernisations in Ukraine
- T-64BM2, with a 57DFM 850-hp (625 kW) engine, a new 1A43U fire control, a new 6ETs43 loader and the possibility to fire the 9M119 missile (NATO code "AT-11 Sniper").
- T-64U which integrated on top the 1A45 fire control (from the T-80U and T-84), PNK-4SU and TKN-4S optics for the tank commander and PZU-7 for the AA machine gun. The tank leader is then able to drive the tank and to use the gun directly if needed. The tank is also known as BM "Bulat".
The two variants are also protected by Kontakt-5 modular reactive armour, able to resist to kinetic energy projectiles, as opposed to the first models which were efficient only against HEAT shaped-charge ammunition. Those two variants could also be remotorised with the 6TDF 1000 hp (735 kW) engine.
As of October 2010, Ukrainian Army has 47 T-64BM "Bulat" [Т-64БМ "Булат"] in service. In 2010 the Kharkiv Malyshev Factory upgraded 10 T-64B tanks to T-64BM "Bulat" standard, and a further 19 will be delivered in 2011. These 29 tanks are being upgraded under a 200 million Hryvnia ($25.1 million) contract signed in April 2009. The T-64B [Т-64Б] tanks were originally produced at Kharkiv in 1980. According to Constantin Isyak (chief engineer of Malyshev Factory), the T-64BM "Bulat" is armoured to the level of modern tanks. They have 'Knife' [Нiж] reactive armour, and the 'Warta' [Варта] active defence system. The T-64BM "Bulat" weighs 45 tonnes (44 long tons), and with its 850 hp (630 kW) 5TDFM multi-fuel diesel engine can do 70 km/h (43 mph), and has a range of 385 km (239 mi). It retains the 125 mm smoothbore gun with an autoloader for 28 rounds, some of which can be guided missiles. It has a 12.7 mm AA machinegun, and a 7.62 mm coaxial machinegun.
Production history
The T-64 first entered production in 1967, shortly before the T-72. (Serial production begin in 1963. The T-64 formally entered service in army in 1967.) The T-64 was KMDB's high-technology offering, intended to replace the IS-3 and T-10 heavy tanks in independent tank battalions. Meanwhile, the T-72 was intended to supersede the T-55 and T-62 in equipping the bulk of Soviet tank and mechanized forces, and for export partners and east-bloc satellite states.
It introduced a new autoloader, which is still used on all T-64s currently in service, as well as all variants of the T-80 except the Ukrainian T-84-120. The T-64 prototypes had the same 115 mm smoothbore gun as the T-62, the ones put in full-scale production had the 125 mm gun.
While the T-64 was the superior tank, it was more expensive and physically complex, and was produced in smaller numbers. The T-72 is mechanically simpler and easier to service in the field, while it is not as well protected, and its manufacturing process is correspondingly simpler. In light of Soviet doctrine, the relatively small numbers of superior T-64 were kept ready and reserved for the most important mission: a potential outbreak of a war in Europe.
The T-64 was never common in Soviet service, except with those units stationed in East Germany. No T-64s were exported. Many T-64s ended up in Russian and Ukrainian service after the breakup of the Soviet Union.
- Ob'yekt 430 (1957) – Prototype with D-10T 100-mm gun, 120 mm armour, 4TPD 580 hp (427 kW) engine, 36 tonnes.
- Ob'yekt 430U – Project, equipped with a 122-mm gun and 160 mm of armour.
- T-64 or Ob'yekt 432 (1961) – Prototype with a D-68 115-mm gun, then initial production version with the same features, about 600 tanks produced.
- T-64R (remontirniy, rebuilt) or Ob'yekt 432R – Redesigned between 1977 and 1981 with external gear from the T-64A but still with the 115-mm gun.
- T-64A or Ob'yekt 434 – 125-mm gun, “gill” armour skirts, a modified sight, and suspension on the fourth road wheel.
- T-64T (1963) – Experimental version with a GTD-3TL 700 hp (515 kW) gas turbine.
- Ob'yekt 436 – Alternative version of Ob'yekt 432 with a V-45 engine. Three built.
- Ob'yekt 438 and Ob'yekt 439 – Ob'yekt 434 with V-45 diesel engine.
- T-64AK or Ob'yekt 446 (1972) – Command version, with a R-130M radio and its 10 m (33 ft) telescoping antenna, a TNA-3 navigation system, without antiarcraft machine gun, carrying 38 rounds of main gun ammunition.
- Ob'yekt 447 – Prototype of the T-64B. Basically a T-64A fitted with the 9K112 "Kobra" system and a1G21 gunsight . This is the "T-64A" displayed in the Kiev museum.
- T-64B or Ob'yekt 447A (1976) – Fitted with redesigned armour, 1A33 fire control system, 9K112-1 "Kobra" ATGM system (NATO code "AT-8 Songster"), TPN-1-49-23 sight, 2A46-2 gun, 2E26M stabiliser and 6ETs40 loader. Later B/BV models have more modern systems 1A33-1, TPN-3-49, 2E42 and a 2A46M-1 gun. From 1985 the T-64B was fitted with stronger glacis armour; older tanks were upgraded with a 16-mm armour plate. Tanks, equipped with the 1,000 hp 6DT engine are known as T-64BM.
- T-64BV – Features "Kontakt-1" reactive armour and "Tucha" 81-mm smoke grenade launchers on the left of the turret.
- T-64BM2 or Ob'yekt 447AM-2 – "Kontakt-5" reactive armour, rubber protection skirts, 1A43U fire control, 6ETs43 loader and able to fire the 9K119 missile (NATO code "AT-11A Sniper"), 5TDFM 850 hp (625 kW) engine.
- T-64U, BM Bulat, or Ob'yekt 447AM-1 – Ukrainian modernisation, bringing the T-64B to the standard of the T-84. Fitted with "Nozh" reactive armour, 9K120 "Refleks" missile (NATO code "AT-11 Sniper"), 1A45 "Irtysh" fire control, TKN-4S commander's sight, PZU-7 antiaircraft machine-gun sight, TPN-4E "Buran-E" night vision, 6TDF 1,000-hp (735 kW) engine.
- T-64B1 or Ob'yekt 437 – Same as the B without the fire control system, carrying 37 shells.
- T-64B1M – T-64Ba equipped with the 1,000-hp 6DT engine.
- T-64BK and T-64B1K or Ob'yekt 446B – Command versions, with an R-130M radio and its 10-m telescoping antenna, a TNA-3 navigation system and AB-1P/30 APU, without antiaircraft machine gun, carrying 28 shells.
- Obyekt 476 – Five prototypes with the 6TDF engine, prototypes for T-80UD development.
- BREM-64 or Ob'yekt 447T – Armoured recovery vehicle with a light 2.5-tonne crane, dozer blade, tow bars, welding equipment, etc. Only a small number was built.
- T-55-64 – Heavily upgraded T-55 with the complete hull and chassis of the T-64, fitted with "Kontakt-1" ERA. Prototype.
- T-80 and T-84 –further developments of the T-64.
- 1977–1981 – brought to the T-64R standard, reorganisation of external equipment as on the T-64A.
- 1972 redesign, fire control improvement (TPD-2-49 and TPN-1-49-23), inclusion of the NSVT machine gun on an electrical turret, R-123M radio.
- 1975 redesign, new 2E28M stabiliser, 6ETs10M loader, multi-fuel engine, 2A46-1 gun and TNPA-65 night vision.
- 1981 redesign, two sets of six 902A smoke grenade launchers, rubber skirts on the suspension instead of the Gill protection.
- 1983 T-64AM,T-64AKM, some tanks were equipped with the 6TDF engine during maintenance.
- 1981 redesign, 2 sets of four 902B2 smoke grenade launchers, 2A26M1 gun.
- 1983 T-64BM,T-64B1M,T-64BMK and T-64B1MK: some tanks were equipped with the 6TDF engine during maintenance.
- 1985 T-64BV,T-64B1V,T-64BVK and T-64B1VK: with "Kontakt" reactive armour, smoke grenade launchers on the left of the turret.
- BM Bulat – T-64 modernization by the Malyshev Factory in Ukraine (see above).
- BMPV-64 – Heavy infantry fighting vehicle, based on the chassis of the T-64 but with a completely redesigned hull with a single entry hatch in the rear. Armament consists of a remote-controlled 30-mm gun. Combat weight is 34.5 tons. The first prototype was ready in 2005.
- BTRV-64 – Similar APC version.
- UMBP-64 – Modified version that will serve as the basis for several (planned) specialized vehicles, including a fire support vehicle, an ambulance and an air-defence vehicle.
- BMPT-K-64 – This variant is not tracked but has a new suspension with 4 axles, similar to the Soviet BTR series. The vehicle is powered by a 5TDF-A/700 engine and has a combat weight of 17.7 tons. It is fitted with a RCWS and can transport 3+8 men. Prototype only.
- BAT-2 – Fast combat engineering vehicle with the engine, lower hull and "small roadwheels" suspension of the T-64. The 40-ton tractor sports a very large, all axis adjustable V-shaped hydraulic dozer blade at the front, a single soil ripper spike at the rear and a 2-ton crane on the top. The crew compartment holds 8 persons (driver, commander, radio operators plus a five-man sapper squad for dismounted tasks). The highly capable BAT-2 was designed to replace the old T-54/AT-T based BAT-M, but WARPAC allies received only small numbers due to its high price and the old and new vehicles served alongside during the late Cold War.
Service history
The tank remained secret for a long time, the West often confusing it with the less-evolved T-72 tank. The T-64 was never exported, and has seen only limited combat experience—in the campaigns against Chechen separatists.
According to David Isby the T-64 first entered service in 1967 with the 41st Guards Tank Division in the Kiev Military District, the suggestion being that this was prudent due to the proximity of the division to the factory, and significant teething problems during induction into service that required constant presence of factory support personnel with the division during acceptance and initial crew and service personnel training on the new type.
- Transnistria - T-64BVs are in service in unknown numbers by the Dnestr separatists.
- Russia – Around 100 are in reserve and 4,000 are probably in storage.
- Ukraine – 2,345 were in service as of 1995, 2,277 as of 2000 and 2,215 as of 2005. Currently around 600 are in service, 1500+ in storage and over 90 from those that are in active service are modernized to T64 BM Bulat.
- Uzbekistan – 100 in service as of 2013.
Potential operators
- Peru - T-64s offered by Ukraine will be part of comparative tests done by the Peruvian Army to find a replacement for their aging T-55s. Between 120 and 170 tanks may be aquired. The T-64 is competing against the T-90S, M1A1 Abrams, Leopard 2A4 and A6, and the T-84.
Former operators
- Soviet Union – Passed on to successor states.
T-64BV technical information
Capabilities and Limitations
The T-64 did not share many drawbacks with the T-72, even if it is often confused with it:
- The automatic loader, hydraulic and not electric, is much faster (loading cycle of 6 to 13 seconds) and more reliable, and less sensitive to jolting when running off-road. It also has a "sequence" fire mode which feeds the gun with shells of the same type in less than 5 seconds. It is also able, in the modern versions, to turn backwards to keep a good speed at the end of the loading sequence.
- Driving seems much less exhausting for the crew, thanks to assisted controls and a more flexible suspension. (Perrett 1987:43)
- The ammunition is stowed at the lower point of the turret shaft, minimizing the risks of destruction by self-detonation.
- Protection, remains able to stop some types of modern projectiles.
- The fire control on the B version is very modern.
- The tank commander's cupola provides good vision, the antiaircraft machine gun can be operated from inside the turret; the commander can also control the main gun sight if necessary.
Additionally, the adoption of the autoloader was highly controversial for several reasons:
- Early versions of the autoloader lacked safety features and were dangerous to the tank crews (especially the gunner, who sits nearby): Limbs could be easily caught in the machinery, leading to horrible injuries and deaths. A sleeve unknowingly snagged on one of the autoloader's moving parts could also drag a crewman into the apparatus upon firing. (Perrett 1987:42)
- The turret was poorly configured to allow the human crew to manually load the gun should the autoloader break. In such situations, rate of fire usually slowed to an abysmal one round per minute as the gunner fumbles with the awkward task of working around the broken machine to load the gun. (Perrett 1987:42)
- While having smaller tank crews (three vs. the usual four) is advantageous since more tanks can theoretically be fielded using the same number of soldiers, there are also serious downsides. Tanks require frequent maintenance and refueling, and much of this is physically demanding work that several people must work together to accomplish. Most of the time, these duties are also performed at the end of a long day of operations, when everyone in the tank is exhausted. Having one less crewman for these tasks increases the strain on the remaining three men and increases the frequency of botched or skipped maintenance. This problem worsens if the tank's commander is also an officer who must often perform other duties such as higher-level meetings, leaving only two men to attend to the tank. (Perrett 1987:42-43) All of this means that tanks with three-man crews are more likely to suffer from performance-degrading human exhaustion, and mechanical failures that take longer to fix and that keep the tank from reaching the battlefield. These problems are exacerbated during prolonged time periods of operations.
- The T-64 was criticized for being too mechanically complex, which resulted in a high breakdown rate. Problems were worst with the suspension system, which was of an entirely new and advanced design on the tank. Due to these problems, teams of civilian mechanics from the T-64 factories were "semi-permanent residents" of Soviet tank units early. (Perrett 1987:43-44)
- Length (gun to the front): 9.295 m.
- Length (without the gun): 6.54 m.
- Breadth: 3.6 m.
- Height: 2.17 m.
- Weight: 42.4 t.
- Engine: 5DTF multifuel (diesel, kerosene and petrol) with 5 opposed cylinders, 10 piston, 13.6 L. Developing 700 hp (515 kW) at 2,800 rpm, consumption of 170 to 200 litres per 100 km.
- Transmission: two lateral gearboxes with seven forward and one backward gear.
- Three internal tanks for a 740 litres fuel capacity, two on the mudguards with 140 litres and two droppable 200 litres tanks on the aft end of the chassis.
- max. road speed: 60.5 km/h.
- max off-road speed: 35 km/h.
- power-to-weight ratio: 16.2 hp/t (11.9 kW/t).
- range: 500 km, 700 km with additional tanks.
- ground pressure: 0.9 kgf/cm2 (88 kPa, 12.8 psi).
- able to ford in 1.8 m of water without preparation and 5 m with snorkels.
- crosses a 2.8 m wide trench.
- crosses a 0.8 m high obstacle.
- max. slope 30°.
- 125 mm smoothbore 2A46M-1 gun (D-81TM) with carousel 6ETs40 loader, 28 shots, fire rate 8 shots per minute, 36 embedded shots (8 x 9M112M "Kobra" (NATO code "AT-8 Songster"), 28 shells). Available shells are all fin-stabilised:
- anti-personnel (APERS) version of the 3UOF-36, 3OVF-22, with several perforating abilities.
- armour-piercing shells (APFSDS) 3UBM-17 or 3UBM-19 or older ones with a supplementary charge giving them an initial speed of about 1800 m/s.
- hollow charge shells, 3VUK-25 or 3UBK-21.
- coaxial machine gun 7.62 mm PKT with 1,250 rounds.
- remote-controlled air-defence machine gun 12.7 mm NSVT "Utyos" with 300 rounds.
- 4+4 (T-64B) or 6+6 (T-64A) 81 mm smoke mortars 902B "Tucha-2".
- The 1A33 fire control system, with:
- Radio control of the 9K112 "Kobra" missiles (NATO code "AT-8 Songster") launched from the gun.
- The 2E28M hydraulic stabiliser (vertical range -5°20' to +15°15')
- The gunner day sight 1G42 with embedded laser telemeter.
- The TPN-1-49-23 active IR night sight.
- The L2G IR projector left of the gun for illumination.
- The 1V517 ballistic calculator.
- The 1B11 anemometric gauge.
- The tank commander's cupola is equipped with:
- The PKN-4S combined day and night sight which allows a 360° vision and to fire the main weapons.
- The PZU-6 AA sight.
- The 2Z20 2-axis electrical stabiliser (vertical range -3° to +70°).
- The TPN-3-49 or TPN-4 and TVN-4 night vision for the driver.
- A R-173M radio.
- An NBC protection, with radiation detectors and global compartment overpressure.
- Two snorkels for crossing rivers with a depth up to 5 m.
- A KMT-6 mine clearing plough can be fitted at the front.
- 3-layer composite armour (K formula), with a thickness between 450 and 20 mm:
- front: 120 mm steel, 105 mm glass fibre, 40 mm steel.
- sides: 80 mm steel.
- front of the turret: 150 mm steel, 150 mm glass fibre, 40 mm steel
- lateral rubber skirts protecting the top of the suspension.
- Kontakt-1 reactive bricks covering:
- the front and the side of the turret
- the glacis
- the lateral skirts
See also
|Wikimedia Commons has media related to: T-64 tanks|
Tanks of comparable role, performance and era
- Chieftain tank : Approximate British equivalent
- T-64A Main Battle Tank at KMDB.
- Три танкиста (Three Tankers)
- Perrett 1987:42
- http://www.meshwar.vistcom.ru/tech/t-64.htm Main battle tank T-64 (Основной боевой танк Т-64)
- wknews.ru Украинская армия получила десять модернизированных Т-64, 28 October 2010
- Kharkiv Morozov Machine Building Design Bureau Main Characteristics of the Upgraded BM Bulat Battle Tank
- Sewell, Stephen, CW2 (rtd). "Why Three Tanks?" (Armor, July–August 1998), p.45.
- Т-64: Чи піде «під ніж» унікальна техніка? (T-64: Will Unique Technology go "Under the Knife"?) at Військо України (Ukrainian Army)
- p.13, Isby, per "Victor Suvorov"
- T-64 MBT at Warfare.ru
- Ground Forces Equipment - Ukraine
- Uzbek-Army Equipment
- Peruvian Tank Contenders - Army-Technology.com, May 17, 2013
- Isby, D.C. (1988). Ten million bayonets: inside the armies of the Soviet Union, Arms and Armour Press, London. ISBN 978-0-85368-774-0
- Perrett, Bryan (1987). Soviet Armour Since 1945. London: Blandford Press. ISBN 0-7137-1735-1.
- Saenko, M., V. Chobitok (2002). Osnovnoj boevoj tank T-64, Moscow: Eksprint. ISBN 5-94038-022-0.
- Sewell, Stephen ‘Cookie’ (1998). “Why Three Tanks?” in Armor vol. 108, no. 4, p. 21. Fort Knox, KY: US Army Armor Center. ISSN 0004-2420. (PDF format)
- Zaloga, Steven (1992), T-64 and T-80, Hong Kong: Concord, ISBN 962-361-031-9.
- BM Bulat Main Battle Tank, Ukraine
- Ukrspets on T-64 upgrades
- Kampfpanzer T-64 (German language)
- T-64 and Bulat at KMDB (manufacturer's site) |
CHRISTIANS have always been interested in God's ancient promise of a
'new covenant.' A covenant is an agreement between two parties. These
questions call attention to points that touch the subject, and help
disciples of Jesus understand New Testament teachings of it.
- 1. Why did Jesus mention the new covenant as being 'in' His blood when
He requested the disciples to drink of the cup in remembrance of Him?
- The words of Jesus were: "This cup is the new covenant in My
blood, which is shed for you"; "drink ye all of it"; "this do ye...in
remembrance of me"; "and they all drank of it" (Luke 22:20, Matt.
26:27, 1 Cor. 11:26, Mark 14:23). Our Master said the new covenant was
'in' His blood. That covenant was dependent upon His death. The
instruction that they drink of the cup meant they were to partake of
the blessings and privileges which that covenant promised. It was a
ceremonial act to convey a truth. Every observance is to be in
remembrance of the gift of His life, upon which all our hopes of
Jesus had a full understanding of all of God's arrangements with
Israel. That understanding brought His teaching that "this cup is the
new covenant in my blood." He knew that an old covenant had been put
in force after animals had been slain, and that the blood of those
animals was used in a prescribed manner (Exod. 24:5-11). He also
understood that God's promise through Jeremiah of a new covenant (Jer.
31:34) made the former covenant old. This fact was later plainly
stated in Hebrews. "When He [God] says 'new' He has made old
the first; and the thing being made old and growing aged is near
vanishing."--Heb. 8:13 Marshall Interlinear
But there was even more intended by our Master's words following
that eventful supper. God had said that through the new covenant He
would "forgive their iniquity, and I will remember their sin no more."
And Jesus knew that the words, "He was wounded for our transgressions,
He was bruised for our iniquities:...for He shall bear their
iniquities" (Isa. 53:5,11) prophesied of Himself and were written for
His instruction as well as for the instruction of all who have faith.
He was to be the victim by which the new covenant would be firm; that
covenant would be dedicated, inaugurated, ratified, put into effect as
a result of His sacrifice (Heb. 9:16-18). The blessing of sonship and
fellowship with God was about to reach others. Jesus would have all
believers acknowledge and remember His sacrifice, its purpose, and the
new arrangement through which their new life and spiritual
relationship with God become a reality.
- 2. Of whom did David prophesy in Psalm 110:4, "Thou art a priest for
ever after the order of Melchizedek," and when did He become that
- Heb. 6:20 answers the first part of the question: "Jesus, made
["who has become" or "having become"] an high priest for ever after
the order of Melchizedec." The quotations in Heb. 5:5,6 and Acts 13:33
from Psalm 2:7 and 110:4 indicate that Christ Jesus became such a
priest at His resurrection. "God hath fulfilled the same unto us their
children, in that He hath raised up Jesus; as it is also written in
the second psalm, Thou art My Son, this day have I begotten Thee." "As
he saith also in another place, Thou art a priest for ever after the
order of Melchizedec." "He became the author [the cause or the source]
of eternal salvation unto all them that obey Him" (Heb. 5:9,10). He
has been made so "after the power of an endless life" in fulfillment
of God's oath given centuries before.--Heb. 7:16-21
- 3. What did Jesus provide which qualified Him to become the
- Jesus is a priest for ever because His sacrifice provided the
blood which sealed the new covenant. The several lines of logic
written regarding Melchizedec in Hebrews 7:5-21 prove that even before
the law established the Levitical priesthood, God planned to terminate
what He foresaw would be an ineffectual ministry. The first reference
in Hebrews to the new covenant is made immediately following the last
reference to Melchizedec in that epistle: "By so much ["to that extent
also" or "because of this oath"] was Jesus made a surety of ["Jesus
has become the guarantee of" or "Jesus has become the one given in
pledge of"] a better covenant" (Heb. 7:22). That arrangement of the
epistle, in which its last mention of Melchizedec was immediately
followed by its first mention of the new covenant, is as though the
writer anticipated that a recognition of Jesus as priest in the higher
order would lead them to discern the reality of the new covenant and
Furthermore, this is suggested by a summary of the chief points
in the first seven chapters of Hebrews. The writer affirms that Jesus
is "a minister in the sanctuary, and in the true tabernacle, which the
Lord pitched, not man." In the type, the Levitical priest ministered
on earth; in the antitype, Jesus ministers in heaven. Regarding the
ministry of Jesus, it is written: "But now He has obtained a more
excellent ministry, by so much as He is mediator of a better covenant
[the new covenant], which has been enacted on better promises."--Heb.
- 4. The blood of what covenant is referred to in Heb. 10:29?
- The reference to blood in vs. 29 is to that of the new covenant.
"...the blood of the covenant" expression in Heb. 9:20 refers back to
the blood of beasts similarly mentioned in Exod. 24:8; "...the blood
of the covenant" in Heb. 10:29 refers to the blood of Jesus. Vs. 29
contrasts the punishment of transgressors in the old covenant with
punishment of such in the new. Both covenants were put in force by
blood, but blood of such merit as to bring eternal redemption--the
blood of Jesus-- is identified only with the new covenant. The "sorer
punishment" would come upon one who "was sanctified" by "the blood of
the covenant" but who came to consider it as common or unholy, as not
alone the only blood efficacious to remit sins. Note that there would
have been no warning to the Hebrews (vs. 25-35) of punishment for
disrespect of "the blood of the covenant" if the covenant ratified by
His blood were not in operation.
Preceding verses are very significant. "For by a single offering
He has perfected for all time those who are sanctified. And the holy
spirit also bears witness to us; for after saying, 'This is the
covenant I will make with them after those days, says the Lord: I will
put my laws on their hearts, and write them on their minds,' then he
adds, 'I will remember their sins and their misdeeds no more.' Where
there is forgiveness of these, there is no longer any offering for
sin" (Heb. 10:14-18 RSV). The holy spirit which communicated prophecy
to Jeremiah also testifies to those who are having God's law put on
their hearts, and written in their minds. It testifies that the
entirely effective and sufficient, and by no means common or unholy
offering for sins made by Jesus, who offered the "one sacrifice for
sins," is that which sanctifies. The apostle's words--"the blood of
the covenant"--relate to those of Jesus: "this cup is the new covenant
in my blood."--Luke 22:20
- 5. God promised to make the new covenant "with the house of Israel"
(Jer. 31:33). Why, then, did the Hebrews' writer quote that prophecy
and write of the new covenant that "a better covenant...has been
enacted" (Heb. 8:6,10 NAS), putting its enactment in past time?
- The inspired author wrote of the new covenant, that it "has been
enacted," because he accepted the witness of the holy spirit. He
recognized the meaning of Jeremiah's prophecy. God had indeed enacted
the new covenant "with the house of Israel." For about three years it
was made only with members of that house, only with Jews. God worded
the promise that way because He foresaw that respective remnants of
the house of Israel and the house of Judah would be together following
their captivities (Ezek. 37:18-22). People of those two houses were
together throughout and following the ministry of Jesus; and He
addressed them as the "house of Israel" (see Matt. 10:6, 15:24).
However, only a relatively few of that 'house' accepted Jesus as
Messiah so as to be received into the covenant. But that is no cause
to conclude that God did not keep His promise to make it with that
house. The fulfillment of his promise is affirmed: "As many as
received him [Jesus], to them gave He power to become the sons of God,
even to them that believe on His name [the name of Christ Jesus our
Cornelius and his group were the first Gentiles received into the
new covenant (Acts 10). They too believed into Jesus. Paul shows that
Gentiles "wert graffed in among them [the natural olive branches], and
with them partakest of the root and fatness of the olive tree" (Rom.
11:17). Believing Gentiles become of the one true Israel (Rom. 11:26),
and are counted among the 12,000 numbered in each of the 12 tribes.--
- 6. To what period does "these days" refer (Jer. 31:33), after which
the new covenant would be made?
- The period, "those days," extended from the beginning of Israel's
desert wanderings through the ministry of Jesus. "Those days" began
when God found fault with Israel (Heb. 8:8). Their acts of
disobedience in the wilderness were so grievous and frequent as to be
recalled by the simple expression, "as in the provocation" (Heb.
3:8,15). See Exod. 16:2, 17:2, 32:1-29, Num. 11:1,6-33, 13:1-14:39,
Ezek. 20:14,22,44. Because the first covenant was not faultless, place
was "sought for the second." Israel frequently "continued not in My
covenant" during the many centuries before Jeremiah prophesied of the
new covenant, and God therefore "regarded them not" (Heb. 8:7-9). But
He kept His covenant, and visited them with necessary adversity,
affliction, and punishment (Amos 3:2). The years after Jeremiah were
also part of "those days" mentioned by the prophet. When the meaning
of the prophecy is misunderstood, an opening is given to the mistaken
idea that "those days" refer to the years during which the gospel has
been preached. And that precludes perception that the new covenant
began to bless believers at Pentecost.
- 7. Is Jesus alone the mediator of the new covenant?
- Yes. There is no statement in Scripture that any but Jesus is
mediator of the new covenant. The apostle Paul indicates clearly who
is mediator. "For there is one God, also one mediator of God and of
man, a man Christ Jesus, the one having given Himself a ransom on
behalf of all, the testimony in its own times" (1 Tim. 2:5,6 Marshall
Interlinear). It was the giving of Himself as ransom for all mankind
that enables Him to be mediator of God and of men. He alone was the
ransom; He alone is mediator. No other person or entity is mentioned
in Scripture as sharing that office with Jesus.
Inasmuch as the better covenant "has been enacted" (Heb. 8:6
NAS), it is evident that Christ Jesus functions as its mediator, for
without a mediator there could be no new covenant and it could not
have been enacted. The holy spirit's testimony to us regarding the
writing of God's law in our hearts (Heb. 10:15-18), which writing
affirms the remission of our sins, is further assurance that Jesus is
mediator of the new covenant.
Moses alone was mediator of the old covenant (Gal. 3:19), and, as
interpreted by the Apostle Peter soon after the day of Pentecost,
Jesus Christ is the greater than Moses prophesied of in Deut.
18:15,18. He quoted that prophecy to the Jews in Jerusalem, together
with its warning that "every soul, which would not hear [so as to
believe and obey] that prophet would be destroyed [Greek, #1842,
exolothreuo; Strong, "to extirpate," a word which Webster defines as
"to pluck up by the stem or root; to eradicate"] from the people [of
God, because not worthy to enjoy fellowship with God's true Israel]."
The Apostle affirmed that all prophets from Samuel on who had spoken,
had all "foretold of these days," the days of Peter's time; and that
"God, having raised up His servant, sent Him to you first, to bless
you in turning every one of you from your iniquities."--Acts 3:22-26
- 8. What blessings do those in the new covenant enjoy?
- Those in relationship with God in the new covenant have peace
with God, sanctification, inner joy, and all other spiritual blessings
as new creatures in Christ Jesus (2 Cor. 5:15,16). Such have come to
God through the name and power of Jesus. "But you have come...to Jesus
the mediator of a new covenant and to the sprinkled blood, which
speaks better than the blood of Abel" (Heb. 12:22-24 NAS). This
reference to the antitypical covenant and its ratification is drawn
from the typical arrangement. After Israel heard and accepted "every
commandment of the law," it was "the blood of the covenant"--animal
blood literally sprinkled on "all the people"--that brought them in
covenant with God at Sinai under Moses, its mediator (Exod. 24:8, Heb.
9:19,20). But there is no need that blood be literally sprinkled upon
those who spiritually "eat the flesh of the Son of man, and drink His
blood," and who receive life thereby; who hear instruction and
endeavor to obey God's will (John 6:53). Their hearts are sprinkled
from a consciousness of evil as they trust and obey.--1 Pet. 1:2, Heb.
Hebrews 9:14,15 verify a blessing received by Jews who had
already come into the new covenant. "How much more, then, will the
blood of Christ, who through the eternal spirit offered Himself
unblemished to God, cleanse our consciences from acts that lead to
death, so that we may serve the living God! For this reason Christ is
the mediator of a new covenant, that those who are called may receive
the promised eternal inheritance--now that He has died as a ransom to
set them free from the sins committed under the first covenant."
Certain elements in the Greek text not disclosed by that NIV rendering
are seen in the Marshall Interlinear translation: "By how much more
the blood of Christ, who through the eternal spirit offered Himself
unblemished to God, will cleanse the conscience of us from dead works
to serve the living God. And therefore He is mediator of a new
covenant, so as death [His death] having occurred for redemption of
the transgressions under the first covenant, the ones having been
called may receive the promise of the eternal inheritance." That
blessing is also for Gentiles, none of whom transgressed the first
covenant because never having been in it, but all of whom nonetheless
missed the mark of keeping God's law.
- 9. Why does Heb. 12:24 use a different Greek word for 'new' than used
elsewhere in the New Testament and in the Septuagint (Greek O.T.) for
- To indicate that the new covenant was then in effect! The Greek
kainos, Strong #2537--"that which is unaccustomed or unused,...new as
to form or quality, of different nature from what is contrasted as
old" (Vine)--was written of the new covenant when prophecy of it was
made, quoted, or referred to (Mark 14:24, Luke 22:20, 1 Cor. 11:25, 2
Cor. 3:6, Heb. 8:8,13, 9:15, Jer. 31:31). God promised an arrangement
"of different nature" than the covenant which had been made old and
was vanishing away. Kainos denoted the difference between the old and
new--a new covenant, of different nature than the old.
But when the writer in Heb. 12:22-24 (NAS) encouraged their
spiritual relationships, the promise had become reality. Thus he wrote
that "you have come...to Jesus the mediator of a new covenant" [Greek
neos, Strong #3501, "new in respect of time, that which is recent"--
Vine's]. It is fitting, in this last appearance in the Bible of the
words 'new covenant,' that those believers were instructed that they
had come to a new, recently-made covenant for forgiveness of sins,
mercy to their unrighteousness, and everlasting life.
- 10. Do disciples have part in the new covenant other than being
blessed in it?
- Yes, indeed! Though they cannot be its mediator, they can surely
be ministers of the new covenant. The glory of its reality is due to
God who planned it, and to Christ whose blood has ratified the new
covenant. Paul said of believers that it is "being manifested that ye
are an epistle of Christ, ministered by us, having been inscribed not
by ink but by the spirit of a living God, not in stony tables but in
tables which are fleshy hearts...the competence of us is God, who also
made us competent as ministers of a new covenant, not of letter but of
spirit, for the letter kills, but the spirit makes alive" (2 Cor.
3:3,6). "And all things are of God, the one having reconciled us to
Himself through Christ and having given to us the ministry of
reconciliation, as that God was in Christ reconciling the world to
Himself, not reckoning to them the trespasses of them, and placing in
us the word of reconciliation. On behalf of Christ therefore we are
ambassadors as God beseeching through us; we beg on behalf of Christ,
Be ye reconciled to God."--2 Cor. 5:18-20 Marshall |
Henry Gray (18211865). Anatomy of the Human Body. 1918.
2e. The Abdomen
The abdomen is the largest cavity in the body. It is of an oval shape, the extremities of the oval being directed upward and downward. The upper extremity is formed by the diaphragm which extends as a dome over the abdomen, so that the cavity extends high into the bony thorax, reaching on the right side, in the mammary line, to the upper border of the fifth rib; on the left side it falls below this level by about 2.5 cm. The lower extremity is formed by the structures which clothe the inner surface of the bony pelvis, principally the Levator ani and Coccygeus on either side. These muscles are sometimes termed the diaphragm of the pelvis. The cavity is wider above than below, and measures more in the vertical than in the transverse diameter. In order to facilitate description, it is artificially divided into two parts: an upper and larger part, the abdomen proper; and a lower and smaller part, the pelvis. These two cavities are not separated from each other, but the limit between them is marked by the superior aperture of the lesser pelvis.
The abdomen proper differs from the other great cavities of the body in being bounded for the most part by muscles and fasciæ, so that it can vary in capacity and shape according to the condition of the viscera which it contains; but, in addition to this, the abdomen varies in form and extent with age and sex. In the adult male, with moderate distension of the viscera, it is oval in shape, but at the same time flattened from before backward. In the adult female, with a fully developed pelvis, it is ovoid with the narrower pole upward, and in young children it is also ovoid but with the narrower pole downward.
Boundaries.It is bounded in front and at the sides by the abdominal muscles and the Iliacus muscles; behind by the vertebral column and the Psoas and Quadratus lumborum muscles; above by the diaphragm; below by the plane of the superior aperture of the lesser pelvis. The muscles forming the boundaries of the cavity are lined upon their inner surfaces by a layer of fascia.
The abdomen contains the greater part of the digestive tube; some of the accessory organs to digestion, viz., the liver and pancreas; the spleen, the kidneys, and the suprarenal glands. Most of these structures, as well as the wall of the cavity in which they are contained, are more or less covered by an extensive and complicated serous membrane, the peritoneum.
The Apertures in the Walls of the Abdomen.The apertures in the walls of the abdomen, for the transmission of structures to or from it, are, in front, the umbilical (in the fetus), for the transmission of the umbilical vessels, the allantois, and vitelline duct; above, the vena caval opening, for the transmission of the inferior vena cava, the aortic hiatus, for the passage of the aorta, azygos vein, and thoracic duct, and the esophageal hiatus, for the esophagus and vagi. Below, there are two apertures on either side: one for the passage of the femoral vessels and lumboinguinal nerve, and the other for the transmission of the spermatic cord in the male, and the round ligament of the uterus in the female.
FIG. 1034 Front view of the thoracic and abdominal viscera. a. Median plane. b b. Lateral planes. c c. Trans tubercular plane. d d. Subcostal plane. e e. Transpyloric plane. (See enlarged image)
Regions.For convenience of description of the viscera, as well as of reference to the morbid conditions of the contained parts, the abdomen is artificially divided into nine regions by imaginary planes, two horizontal and two sagittal, passing through the cavity, the edges of the planes being indicated by lines drawn on the surface of the body. Of the horizontal planes the upper or transpyloric is indicated by a line encircling the body at the level of a point midway between the jugular notch and the symphysis pubis, the lower by a line carried around the trunk at the level of a point midway between the transpyloric and the symphysis pubis. The latter is practically the intertubercular plane of Cunningham, who pointed out163 that its level corresponds with the prominent and easily defined tubercle on the iliac crest about 5 cm. behind the anterior superior iliac spine. By means of these imaginary planes the abdomen is divided into three zones, which are named from above downward the subcostal, umbilical, and hypogastric zones. Each of these is further subdivided into three regions by the two sagittal planes, which are indicated on the surface by lines drawn vertically through points half-way between the anterior superior iliac spines and the symphysis pubis.164
The middle region of the upper zone is called the epigastric; and the two lateral regions, the right and left hypochondriac. The central region of the middle zone is the umbilical; and the two lateral regions, the right and left lumbar. The middle region of the lower zone is the hypogastric or pubic region; and the lateral regions are the right and left iliac or inguinal(Fig. 1034).
The pelvis is that portion of the abdominal cavity which lies below and behind a plane passing through the promontory of the sacrum, lineæ terminales of the hip bones, and the pubic crests. It is bounded behind by the sacrum, coccyx, Piriformes, and the sacrospinous and sacrotuberous ligaments; in front and laterally by the pubes and ischia and Obturatores interni; above it communicates with the abdomen proper; below it is closed by the Levatores ani and Coccygei and the urogenital diaphragm. The pelvis contains the urinary bladder, the sigmoid colon and rectum, a few coils of the small intestine, and some of the generative organs.
When the anterior abdominal wall is removed, the viscera are partly exposed as follows: above and to the right side is the liver, situated chiefly under the shelter of the right ribs and their cartilages, but extending across the middle line and reaching for some distance below the level of the xiphoid process. To the left of the liver is the stomach, from the lower border of which an apron-like fold of peritoneum, the greater omentum, descends for a varying distance, and obscures, to a greater or lesser extent, the other viscera. Below it, however, some of the coils of the small intestine can generally be seen, while in the right and left iliac regions respectively the cecum and the iliac colon are partly exposed. The bladder occupies the anterior part of the pelvis, and, if distended, will project above the symphysis pubis; the rectum lies in the concavity of the sacrum, but is usually obscured by the coils of the small intestine. The sigmoid colon lies between the rectum and the bladder.
When the stomach is followed from left to right it is seen to be continuous with the first part of the small intestine, or duodenum, the point of continuity being marked by a thickened ring which indicates the position of the pyloric valve. The duodenum passes toward the under surface of the liver, and then, curving downward, is lost to sight. If, however, the greater omentum be thrown upward over the chest, the inferior part of the duodenum will be observed passing across the vertebral column toward the left side, where it becomes continuous with the coils of the jejunum and ileum. These measure some 6 meters in length, and if followed downward the ileum will be seen to end in the right iliac fossa by opening into the cecum, the commencement of the large intestine. From the cecum the large intestine takes an arched course, passing at first upward on the right side, then across the middle line and downward on the left side, and forming respectively the ascending transverse, and descending parts of the colon. In the pelvis it assumes the form of a loop, the sigmoid colon, and ends in the rectum.
The glistening appearance of the deep surface of the abdominal wall and of the surfaces of the exposed viscera is due to the fact that the former is lined, and the latter are more or less completely covered, by a serous membrane, the peritoneum.
the Peritoneum (Tunica Serosa)The peritoneum is the largest serous membrane in the body, and consists, in the male, of a closed sac, a part of which is applied against the abdominal parietes, while the remainder is reflected over the contained viscera. In the female the peritoneum is not a closed sac, since the free ends of the uterine tubes open directly into the peritoneal cavity. The part which lines the parietes is named the parietal portion of the peritoneum; that which is reflected over the contained viscera constitutes the visceral portion of the peritoneum. The free surface of the membrane is smooth, covered by a layer of flattened mesothelium, and lubricated by a small quantity of serous fluid. Hence the viscera can glide freely against the wall of the cavity or upon one another with the least possible amount of friction. The attached surface is rough, being connected to the viscera and inner surface of the parietes by means of areolar tissue, termed the subserous areolar tissue. The parietal portion is loosely connected with the fascial lining of the abdomen and pelvis, but is more closely adherent to the under surface of the diaphragm, and also in the middle line of the abdomen.
The space between the parietal and visceral layers of the peritoneum is named the peritoneal cavity; but under normal conditions this cavity is merely a potential one, since the parietal and visceral layers are in contact. The peritoneal cavity gives off a large diverticulum, the omental bursa, which is situated behind the stomach and adjoining structures; the neck of communication between the cavity and the bursa is termed the epiploic foramen (foramen of Winslow). Formerly the main portion of the cavity was described as the greater, and the omental bursa as the lesser sac.
The peritoneum differs from the other serous membranes of the body in presenting a much more complex arrangement, and one that can be clearly understood only by following the changes which take place in the digestive tube during its development.
To trace the membrane from one viscus to another, and from the viscera to the parietes, it is necessary to follow its continuity in the vertical and horizontal directions, and it will be found simpler to describe the main portion of the cavity and the omental bursa separately.
Vertical Disposition of the Main Peritoneal Cavity (greater sac) (Fig. 1035).It is convenient to trace this from the back of the abdominal wall at the level of the umbilicus. On following the peritoneum upward from this level it is seen to be reflected around a fibrous cord, the ligamentum teres (obliterated umbilical vein), which reaches from the umbilicus to the under surface of the liver. This reflection forms a somewhat triangular fold, the falciform ligament of the liver, attaching the upper and anterior surfaces of the liver to the diaphragm and abdominal wall. With the exception of the line of attachment of this ligament the peritoneum covers the whole of the under surface of the anterior part of the diaphragm, and is continued from it on to the upper surface of the right lobe of the liver as the superior layer of the coronary ligament, and on to the upper surface of the left lobe as the superior layer of the left triangular ligament of the liver. Covering the upper and anterior surfaces of the liver, it is continued around its sharp margin on to the under surface, where it presents the following relations: (a) It covers the under surface of the right lobe and is reflected from the back part of this on to the right suprarenal gland and upper extremity of the right kidney, forming in this situation the inferior layer of the coronary ligament; a special fold, the hepatorenal ligament, is frequently present between the inferior surface of the liver and the front of the kidney. From the kidney it is carried downward to the duodenum and right colic flexure and medialward in front of the inferior vena cava, where it is continuous with the posterior wall of the omental bursa. Between the two layers of the coronary ligament there is a large triangular surface of the liver devoid of peritoneal covering; this is named the bare area of the liver, and is attached to the diaphragm by areolar tissue. Toward the right margin of the liver the two layers of the coronary ligament gradually approach each other, and ultimately fuse to form a small triangular fold connecting the right lobe of the liver to the diaphragm, and named the right triangular ligament of the liver. The apex of the triangular bare area corresponds with the point of meeting of the two layers of the coronary ligament, its base with the fossa for the inferior vena cava. (b) It covers the lower surface of the quadrate lobe, the under and lateral surfaces of the gall-bladder, and the under surface and posterior border of the left lobe; it is then reflected from the upper surface of the left lobe to the diaphragm as the inferior layer of the left triangular ligament, and from the porta of the liver and the fossa for the ductus venosus to the lesser curvature of the stomach and the first 2.5 cm. of the duodenum as the anterior layer of the hepatogastric and hepatoduodenal ligaments, which together constitute the lesser omentum. If this layer of the lesser omentum be followed to the right it will be found to turn around the hepatic artery, bile duct, and portal vein, and become continuous with the anterior wall of the omental bursa, forming a free folded edge of peritoneum. Traced downward, it covers the antero-superior surface of the stomach and the commencement of the duodenum, and is carried down into a large free fold, known as the gastrocolic ligament or greater omentum. Reaching the free margin of this fold, it is reflected upward to cover the under and posterior surfaces of the transverse colon, and thence to the posterior abdominal wall as the inferior layer of the transverse mesocolon. It reaches the abdominal wall at the head and anterior border of the pancreas, is then carried down over the lower part of the head and over the inferior surface of the pancreas on the superior mesenteric vessels, and thence to the small intestine as the anterior layer of the mesentery. It encircles the intestine, and subsequently may be traced, as the posterior layer of the mesentery, upward and backward to the abdominal wall. From this it sweeps down over the aorta into the pelvis, where it invests the sigmoid colon, its reduplication forming the sigmoid mesocolon. Leaving first the sides and then the front of the rectum, it is reflected on to the seminal vesicles and fundus of the urinary bladder and, after covering the upper surface of that viscus, is carried along the medial and lateral umbilical ligaments (Fig. 1036) on to the back of the abdominal wall to the level from which a start was made.
FIG. 1035 Vertical disposition of the peritoneum. Main cavity, red; omental bursa, blue. (See enlarged image)
FIG. 1036 Posterior view of the anterior abdominal wall in its lower half. The peritoneum is in place, and the various cords are shining through. (After Joessel.) (See enlarged image)
Between the rectum and the bladder it forms, in the male, a pouch, the rectovesical excavation, the bottom of which is slightly below the level of the upper ends of the vesiculæ seminalesi. e., about 7.5 cm. from the orifice of the anus. When the bladder is distended, the peritoneum is carried up with the expanded viscus so that a considerable part of the anterior surface of the latter lies directly against the abdominal wall without the intervention of peritoneal membrane (prevesical space of Retzius). In the female the peritoneum is reflected from the rectum over the posterior vaginal fornix to the cervix and body of the uterus, forming the rectouterine excavation (pouch of Douglas). It is continued over the intestinal surface and fundus of the uterus on to its vesical surface, which it covers as far as the junction of the body and cervix uteri, and then to the bladder, forming here a second, but shallower, pouch, the vesicouterine excavation. It is also reflected from the sides of the uterus to the lateral walls of the pelvis as two expanded folds, the broad ligaments of the uterus, in the free margin of each of which is the uterine tube.
Vertical Disposition of the Omental Bursa (lesser peritoneal sac) (Fig. 1035).A start may be made in this case on the posterior abdominal wall at the anterior border of the pancreas. From this region the peritoneum may be followed upward over the pancreas on to the inferior surface of the diaphragm, and thence on to the caudate lobe and caudate process of the liver to the fossa from the ductus venosus and the porta of the liver. Traced to the right, it is continuous over the inferior vena cava with the posterior wall of the main cavity. From the liver it is carried downward to the lesser curvature of the stomach and the commencement of the duodenum as the posterior layer of the lesser omentum, and is continuous on the right, around the hepatic artery, bile duct, and portal vein, with the anterior layer of this omentum. The posterior layer of the lesser omentum is carried down as a covering for the postero-inferior surfaces of the stomach and commencement of the duodenum, and is continued downward as the deep layer of the gastrocolic ligament or greater omentum. From the free margin of this fold it is reflected upward on itself to the anterior and superior surfaces of the transverse colon, and thence as the superior layer of the transverse mesocolon to the anterior border of the pancreas, the level from which a start was made. It will be seen that the loop formed by the wall of the omental bursa below the transverse colon follows, and is closely applied to, the deep surface of that formed by the peritoneum of the main cavity, and that the greater omentum or large fold of peritoneum which hangs in front of the small intestine therefore consists of four layers, two anterior and two posterior separated by the potential cavity of the omental bursa.
Horizontal Disposition of the Peritoneum.Below the transverse colon the arrangement is simple, as it includes only the main cavity; above the level of the transverse colon it is more complicated on account of the existence of the omental bursa. Below the transverse colon it may be considered in the two regions, viz., in the pelvis and in the abdomen proper.
FIG. 1037 The peritoneum of the male pelvis. (Dixon and Birmingham.) (See enlarged image)
(1) In the Pelvis.The peritoneum here follows closely the surfaces of the pelvic viscera and the inequalities of the pelvic walls, and presents important differences in the two sexes. (a) In the male(Fig. 1037) it encircles the sigmoid colon, from which it is reflected to the posterior wall of the pelvis as a fold, the sigmoid mesocolon. It then leaves the sides and, finally, the front of the rectum, and is continued on to the upper ends of the seminal vesicles and the bladder; on either side of the rectum it forms a fossa, the pararectal fossa, which varies in size with the distension of the rectum. In front of the rectum the peritoneum forms the rectovesical excavation, which is limited laterally by peritoneal folds extending from the sides of the bladder to the rectum and sacrum. These folds are known from their position as the rectovesical or sacrogenital folds. The peritoneum of the anterior pelvic wall covers the superior surface of the bladder, and on either side of this viscus forms a depression, termed the paravesical fossa, which is limited laterally by the fold of peritoneum covering the ductus deferens. The size of this fossa is dependent on the state of distension of the bladder; when the bladder is empty, a variable fold of peritoneum, the plica vesicalis transversa, divides the fossa into two portions. On the peritoneum between the paravesical and pararectal fossæ the only elevations are those produced by the ureters and the hypogastric vessels. (b) In the female, pararectal and paravesical fossæ similar to those in the male are present: the lateral limit of the paravesical fossa is the peritoneum investing the round ligament of the uterus. The rectovesical excavation is, however, divided by the uterus and vagina into a small anterior vesicouterine and a large, deep, posterior rectouterine excavation. The sacrogenital folds form the margins of the latter, and are continued on to the back of the uterus to form a transverse fold, the torus uterinus. The broad ligaments extend from the sides of the uterus to the lateral walls of the pelvis; they contain in their free margins the uterine tubes, and in their posterior layers the ovaries. Below, the broad ligaments are continuous with the peritoneum on the lateral walls of the pelvis. On the lateral pelvic wall behind the attachment of the broad ligament, in the angle between the elevations produced by the diverging hypogastric and external iliac vessels is a slight fossa, the ovarian fossa, in which the ovary normally lies.
FIG. 1038 Horizontal disposition of the peritoneum in the lower part of the abdomen. (See enlarged image)
(2) In the Lower Abdomen (Fig. 1038).Starting from the linea alba, below the level of the transverse colon, and tracing the continuity of the peritoneum in a horizontal direction to the right, the membrane covers the inner surface of the abdominal wall almost as far as the lateral border of the Quadratus lumborum; it encloses the cecum and vermiform process, and is reflected over the sides and front of the ascending colon; it may then be traced over the duodenum, Psoas major, and inferior vena cava toward the middle line, whence it passes along the mesenteric vessels to invest the small intestine, and back again to the large vessels in front of the vertebral column, forming the mesentery, between the layers of which are contained the mesenteric bloodvessels, lacteals, and glands. It is then continued over the left Psoas; it covers the sides and front of the descending colon, and, reaching the abdominal wall, is carried on it to the middle line.
FIG. 1039 Horizontal disposition of the peritoneum in the upper part of the abdomen. (See enlarged image)
(3) In the Upper Abdomen (Fig. 1039).Above the transverse colon the omental bursa is superadded to the general sac, and the communication of the two cavities with one another through the epiploic foramen can be demonstrated.
(a) Main Cavity.Commencing on the posterior abdominal wall at the inferior vena cava, the peritoneum may be followed to the right over the front of the suprarenal gland and upper part of the right kidney on to the antero-lateral abdominal wall. From the middle line of the anterior wall a backwardly directed fold encircles the obliterated umbilical vein and forms the falciform ligament of the liver. Continuing to the left, the peritoneum lines the antero-lateral abdominal wall and covers the lateral part of the front of the left kidney, and is reflected to the posterior border of the hilus of the spleen as the posterior layer of the phrenicolienal ligament. It can then be traced around the surface of the spleen to the front of the hilus, and thence to the cardiac end of the greater curvature of the stomach as the anterior layer of the gastrolienal ligament. It covers the antero-superior surfaces of the stomach and commencement of the duodenum, and extends up from the lesser curvature of the stomach to the liver as the anterior layer of the lesser omentum.
(b) Omental Bursa (bursa omentalis; lesser peritoneal sac).On the posterior abdominal wall the peritoneum of the general cavity is continuous with that of the omental bursa in front of the inferior vena cava. Starting from here, the bursa may be traced across the aorta and over the medial part of the front of the left kidney and diaphragm to the hilus of the spleen as the anterior layer of the phrenicolienal ligament. From the spleen it is reflected to the stomach as the posterior layer of the gastrosplenic ligament. It covers the postero-inferior surfaces of the stomach and commencement of the duodenum, and extends upward to the liver as the posterior layer of the lesser omentum; the right margin of this layer is continuous around the hepatic artery, bile duct, and portal vein, with the wall of the general cavity.
The epiploic foramen (foramen epiploicum; foramen of Winslow) is the passage of communication between the general cavity and the omental bursa. It is bounded in front by the free border of the lesser omentum, with the common bile duct, hepatic artery, and portal vein between its two layers; behind by the peritoneum covering the inferior vena cava; above by the peritoneum on the caudate process of the liver, and below by the peritoneum covering the commencement of the duodenum and the hepatic artery, the latter passing forward below the foramen before ascending between the two layers of the lesser omentum.
The boundaries of the omental bursa will now be evident. It is bounded in front, from above downward, by the caudate lobe of the liver, the lesser omentum, the stomach, and the anterior two layers of the greater omentum. Behind, it is limited, from below upward, by the two posterior layers of the greater omentum, the transverse colon, and the ascending layer of the transverse mesocolon, the upper surface of the pancreas, the left suprarenal gland, and the upper end of the left kidney. To the right of the esophageal opening of the stomach it is formed by that part of the diaphragm which supports the caudate lobe of the liver. Laterally, the bursa extends from the epiploic foramen to the spleen, where it is limited by the phrenicolienal and gastrolienal ligaments.
The omental bursa, therefore, consists of a series of pouches or recesses to which the following terms are applied: (1) the vestibule, a narrow channel continued from the epiploic foramen, over the head of the pancreas to the gastropancreatic fold; this fold extends from the omental tuberosity of the pancreas to the right side of the fundus of the stomach, and contains the left gastric artery and coronary vein; (2) the superior omental recess, between the caudate lobe of the liver and the diaphragm; (3) the lienal recess, between the spleen and the stomach; (4) the inferior omental recess, which comprises the remainder of the bursa.
In the fetus the bursa reaches as low as the free margin of the greater omentum, but in the adult its vertical extent is usually more limited owing to adhesions between the layers of the omentum. During a considerable part of fetal life the transverse colon is suspended from the posterior abdominal wall by a mesentery of its own, the two posterior layers of the greater omentum passing at this stage in front of the colon. This condition occasionally persists throughout life, but as a rule adhesion occurs between the mesentery of the transverse colon and the posterior layer of the greater omentum, with the result that the colon appears to receive its peritoneal covering by the splitting of the two posterior layers of the latter fold. In the adult the omental bursa intervenes between the stomach and the structures on which that viscus lies, and performs therefore the functions of a serous bursa for the stomach.
Numerous peritoneal folds extend between the various organs or connect them to the parietes; they serve to hold the viscera in position, and, at the same time, enclose the vessels and nerves proceeding to them. They are grouped under the three headings of ligaments, omenta, and mesenteries.
The lesser omentum (omentum minus; small omentum; gastrohepatic omentum) is the duplicature which extends to the liver from the lesser curvature of the stomach and the commencement of the duodenum. It is extremely thin, and is continuous with the two layers of peritoneum which cover respectively the antero-superior and postero-inferior surfaces of the stomach and first part of the duodenum. When these two layers reach the lesser curvature of the stomach and the upper border of the duodenum, they join together and ascend as a double fold to the porta of the liver; to the left of the porta the fold is attached to the bottom of the fossa for the ductus venosus, along which it is carried to the diaphragm, where the two layers separate to embrace the end of the esophagus. At the right border of the omentum the two layers are continuous, and form a free margin which constitutes the anterior boundary of the epiploic foramen. The portion of the lesser omentum extending between the liver and stomach is termed the hepatogastric ligament, while that between the liver and duodenum is the hepatoduodenal ligament. Between the two layers of the lesser omentum, close to the right free margin, are the hepatic artery, the common bile duct, the portal vein, lymphatics, and the hepatic plexus of nervesall these structures being enclosed in a fibrous capsule (Glissons capsule). Between the layers of the lesser omentum, where they are attached to the stomach, run the right and left gastric vessels.
The greater omentum (omentum majus; great omentum; gastrocolic omentum) is the largest peritoneal fold. It consists of a double sheet of peritoneum, folded on itself so that it is made up of four layers. The two layers which descend from the stomach and commencement of the duodenum pass in front of the small intestines, sometimes as low down as the pelvis; they then turn upon themselves, and ascend again as far as the transverse colon, where they separate and enclose that part of the intestine. These individual layers may be easily demonstrated in the young subject, but in the adult they are more or less inseparably blended. The left border of the greater omentum is continuous with the gastrolienal ligament; its right border extends as far as the commencement of the duodenum. The greater omentum is usually thin, presents a cribriform appearance, and always contains some adipose tissue, which in fat people accumulates in considerable quantity. Between its two anterior layers, a short distance from the greater curvature of the stomach, is the anastomosis between the right and left gastroepiploic vessels.
The mesentery proper (mesenterium) is the broad, fan-shaped fold of peritoneum which connects the convolutions of the jejunum and ileum with the posterior wall of the abdomen. Its rootthe part connected with the structures in front of the vertebral columnis narrow, about 15 cm. long, and is directed obliquely from the duodenojejunal flexure at the left side of the second lumbar vertebra to the right sacroiliac articulation (Fig. 1040). Its intestinal border is about 6 metres long; and here the two layers separate to enclose the intestine, and form its peritoneal coat. It is narrow above, but widens rapidly to about 20 cm., and is thrown into numerous plaits or folds. It suspends the small intestine, and contains between its layers the intestinal branches of the superior mesenteric artery, with their accompanying veins and plexuses of nerves, the lacteal vessels, and mesenteric lymph glands.
The transverse mesocolon (mesocolon transversum) is a broad fold, which connects the transverse colon to the posterior wall of the abdomen. It is continuous with the two posterior layers of the greater omentum, which, after separating to surround the transverse colon, join behind it, and are continued backward to the vertebral column, where they diverge in front of the anterior border of the pancreas. This fold contains between its layers the vessels which supply the transverse colon.
The sigmoid mesocolon (mesocolon sigmoideum) is the fold of peritoneum which retains the sigmoid colon in connection with the pelvic wall. Its line of attachment forms a V-shaped curve, the apex of the curve being placed about the point of division of the left common iliac artery. The curve beings on the medial side of the left Psoas major, and runs upward and backward to the apex, from which it bends sharply downward, and ends in the median plane at the level of the third sacral vertebra. The sigmoid and superior hemorrhoidal vessels run between the two layers of this fold.
In most cases the peritoneum covers only the front and sides of the ascending and descending parts of the colon. Sometimes, however, these are surrounded by the serous membrane and attached to the posterior abdominal wall by an ascending and a descending mesocolon respectively. A fold of peritoneum, the phrenicocolic ligament, is continued from the left colic flexure to the diaphragm opposite the tenth and eleventh ribs; it passes below and serves to support the spleen, and therefore has received the name of sustentaculum lienis.
FIG. 1040 Diagram devised by Delépine to show the lines along which the peritoneum leaves the wall of the abdomen to invest the viscera. (See enlarged image)
The appendices epiploicæ are small pouches of the peritoneum filled with fat and situated along the colon and upper part of the rectum. They are chiefly appended to the transverse and sigmoid parts of the colon.
Peritoneal Recesses or Fossæ (retroperitoneal fossæ).In certain parts of the abdominal cavity there are recesses of peritoneum forming culs-de-sac or pouches, which are of surgical interest in connection with the possibility of the occurrence of retroperitoneal herniæ. The largest of these is the omental bursa (already described), but several others, of smaller size, require mention, and may be divided into three groups, viz.: duodenal, cecal, and intersigmoid.
1. Duodenal Fossæ (Figs. 1041,1042).Three are fairly constant, viz.: (a) The inferior duodenal fossa, present in from 70 to 75 per cent. of cases, is situated opposite the third lumbar vertebra on the left side of the ascending portion of the duodenum. Its opening is directed upward, and is bounded by a thin sharp fold of peritoneum with a concave margin, called the duodenomesocolic fold. The tip of the index finger introduced into the fossa under the fold passes some little distance behind the ascending portion of the duodenum. (b) The superior duodenal fossa, present in from 40 to 50 per cent. of cases, often coexists with the inferior one, and its orifice looks downward. It lies on the left of the ascending portion of the duodenum, in front of the second lumbar vertebra, and behind a sickle-shaped fold of peritoneum, the duodenojejunal fold, and has a depth of about 2 cm. (c) The duodenojejunal fossa exists in from 15 to 20 per cent. of cases, but has never yet been found in conjunction with the other forms of duodenal fossæ it can be seen by pulling the jejunum downward and to the right, after the transverse colon has been pulled upward. It is bounded above by the pancreas, to the right by the aorta, and to the left by the kidney; beneath is the left renal vein. It has a depth of from 2 to 3 cm., and its orifice, directed downward and to the right, is nearly circular and will admit the tip of the little finger.
FIG. 1044 Inferior ileocecal fossa. The cecum and ascending colon have been drawn lateralward and downward, the ileum upward and backward, and the vermiform process downward. (Poirier and Charpy.) (See enlarged image)
2. Cecal Fossæ (pericecal folds or fossæ).There are three principal pouches or recesses in the neighborhood of the cecum (Figs. 1043to 1045): (a) The superior ileocecal fossa is formed by a fold of peritoneum, arching over the branch of the ileocolic artery which supplies the ileocolic junction. The fossa is a narrow chink situated between the mesentery of the small intestine, the ileum, and the small portion of the cecum behind. (b) The inferior ileocecal fossa is situated behind the angle of junction of the ileum and cecum. It is formed by the ileocecal fold of peritoneum (bloodless fold of Treves), the upper border of which is fixed to the ileum, opposite its mesenteric attachment, while the lower border, passing over the ileocecal junction, joins the mesenteriole of the vermiform process, and sometimes the process itself. Between this fold and the mesenteriole of the vermiform process is the inferior ileocecal fossa. It is bounded above by the posterior surface of the ileum and the mesentery; in front and below by the ileocecal fold, and behind by the upper part of the mesenteriole of the vermiform process. (c) The cecal fossa is situated immediately behind the cecum, which has to be raised to bring it into view. It varies much in size and extent. In some cases it is sufficiently large to admit the index finger, and extends upward behind the ascending colon in the direction of the kidney; in others it is merely a shallow depression. It is bounded on the right by the cecal fold, which is attached by one edge to the abdominal wall from the lower border of the kidney to the iliac fossa and by the other to the postero-lateral aspect of the colon. In some instances additional fossæ, the retrocecal fossæ, are present.
3. The intersigmoid fossa (recessus intersigmoideus) is constant in the fetus and during infancy, but disappears in a certain percentage of cases as age advances. Upon drawing the sigmoid colon upward, the left surface of the sigmoid mesocolon is exposed, and on it will be seen a funnel-shaped recess of the peritoneum, lying on the external iliac vessels, in the interspace between the Psoas and Iliacus muscles. This is the orifice leading to the intersigmoid fossa, which lies behind the sigmoid mesocolon, and in front of the parietal peritoneum. The fossa varies in size; in some instances it is a mere dimple, whereas in others it will admit the whole of the index finger.165 |
Drugs A - Z
Generic Name: cyanocobalamin | Brand Name: Neuroforte-R
CategoryHerbs & Supplements
B-12, B Complex, B Complex Vitamin, bedumil, cobalamin, cobalamins, cobamin, cyanocobalamin, cyanocobalaminum, cycobemin, hydroxocobalamin, hydroxocobalaminum, hydroxocobemine, idrossocobalamina, methylcobalamin, vitadurin, vitamin B-12.
Vitamin B12 is an essential water-soluble vitamin that is commonly found in a variety of foods such as fish, shellfish, meat, and dairy products. Vitamin B12 is frequently used in combination with other B vitamins in a vitamin B complex formulation. It helps maintain healthy nerve cells and red blood cells and is also needed to make DNA, the genetic material in all cells. Vitamin B12 is bound to the protein in food. Hydrochloric acid in the stomach releases B12 from protein during digestion. Once released, B12 combines with a substance called intrinsic factor (IF) before it is absorbed into the bloodstream.
The human body stores several years' worth of vitamin B12, so nutritional deficiency of this vitamin is extremely rare. Elderly are the most at risk. However, deficiency can result from being unable to use vitamin B12. Inability to absorb vitamin B12 from the intestinal tract can be caused by a disease known as pernicious anemia. Additionally, strict vegetarians or vegans who are not taking in proper amounts of B12 are also prone to a deficiency state.
A day's supply of vitamin B12 can be obtained by eating 1 chicken breast plus 1 hard-boiled egg plus 1 cup plain low-fat yogurt, or 1 cup milk plus 1 cup raisin bran.
EvidenceDISCLAIMER: These uses have been tested in humans or animals. Safety and effectiveness have not always been proven. Some of these conditions are potentially serious, and should be evaluated by a qualified healthcare provider.
Megaloblastic anemia - due to vitamin B12 deficiency:
Vitamin B12 deficiency is a cause of megaloblastic anemia. In this type of anemia, red blood cells are larger than normal and the ratio of nucleus size to cell cytoplasm is increased. There are other potential causes of megaloblastic anemia, including folate deficiency or various inborn metabolic disorders. If the cause is B12 deficiency, then treatment with B12 is the standard approach. Patients with anemia should be evaluated by a physician in order to diagnose and address the underlying cause.
Pernicious anemia (blood abnormality) is a form of anemia that occurs when there is an absence of intrinsic factor, a substance normally present in the stomach. Vitamin B12 binds with intrinsic factor before it is absorbed and used by the body. An absence of intrinsic factor prevents normal absorption of B12 and may result in pernicious anemia. Pernicious anemia treatment is usually lifelong; supplemental vitamin B12 given intramuscularly, intranasally, or by mouth.
Vitamin B12 deficiency:
Studies have shown that a deficiency of vitamin B12 can lead to abnormal neurologic and psychiatric symptoms. These symptoms may include: ataxia (shaky movements and unsteady gait), muscle weakness, spasticity, incontinence, hypotension (low blood pressure), vision problems, dementia, psychoses, and mood disturbances. Researchers report that these symptoms may occur when vitamin B12 levels are just slightly lower than normal and are considerably above the levels normally associated with anemia. People at risk for vitamin B12 deficiency include strict vegetarians, elderly people, breastfed infants, and people with increased vitamin B12 requirements associated with pregnancy, thyrotoxicosis, hemolytic anemia, hemorrhage, malignancy, liver or kidney disease.
Administering vitamin B12 orally, intramuscularly, or intranasally is effective for preventing and treating dietary vitamin B12 deficiency.
Some patients diagnosed with Alzheimer's disease have been found to have abnormally low vitamin B12 levels in their blood. However, vitamin B12 deficiency itself often causes disorientation and confusion and thus mimics some of the prominent symptoms of Alzheimer's disease. Well-designed clinical trials are needed before a strong recommendation can be made.
Some evidence suggests that folic acid plus vitamin B12 and pyridoxine daily can decrease the rate of restenosis in patients treated with balloon angioplasty. But this combination does not seem to be as effective for reducing restenosis in patients after coronary stenting. Due to the lack of evidence of benefit and potential for harm, this combination of vitamins should not be recommended for patients receiving coronary stents.
Researchers at Johns Hopkins University report that women with breast cancer tend to have lower vitamin B12 levels in their blood serum than do women without breast cancer. In a subsequent review of these findings, it was hypothesized that vitamin B12 deficiency may lead to breast cancer because it could result in less folate being available to ensure proper DNA replication and repair. Higher dietary folate intake is associated with a reduced risk of breast cancer. The risk may be further reduced in women who also consume high amounts of dietary vitamin B12 in combination with dietary pyridoxine (vitamin B6) and methionine. However, there is no evidence that dietary vitamin B12 alone reduces the risk of breast cancer.
Hyperhomocysteinemia (high homocysteine levels in the blood) is a risk factor for coronary, cerebral, and peripheral atherosclerosis, recurrent thromboembolism, deep vein thrombosis, myocardial infarction (heart attack), and ischemic stroke. Elevated homocysteine levels may be a marker instead of a cause of vascular disease. However, it is not clear if lowering homocysteine levels results in reduced cardiovascular morbidity and mortality. Folic acid, pyridoxine (vitamin B6), and vitamin B12 supplementation can reduce total homocysteine levels; however, this reduction does not seem to help with secondary prevention of death or cardiovascular events such as stroke or myocardial infarction in people with prior stroke. More evidence is needed to fully explain the association of total homocysteine levels with vascular risk and the potential use of vitamin supplementation.
There is some evidence that intramuscular injections of vitamin B12 given twice per week might improve the general well-being and happiness of patients complaining of tiredness or fatigue. However, fatigue has many potential causes. Well-designed clinical trials are needed before a recommendation can be made.
Some evidence suggests that vitamin B12 in combination with fish oil might be superior to fish oil alone when used daily to reduce total serum cholesterol and triglycerides. Well-designed clinical trials of vitamin B12 supplementation alone are needed before a conclusion can be drawn.
Administering vitamin B12 intramuscularly seems to be effective for treating familial selective vitamin B12 malabsorption (Imerslund-Grasbeck disease). Further research is needed to confirm these results.
Preliminary clinical reports show that cyanocobalamin may help relieve tremor associated with shaky-leg syndrome. Further research is needed to confirm these results.
Sickle cell disease:
One study suggests that a practical daily combination may include folic acid, vitamin B12, and vitamin B6. This combination may be a simple and relatively inexpensive way to reduce these patients' inherently high risk of endothelial damage. Further research is needed to confirm these results.
Circadian rhythm sleep disorders:
Taking vitamin B12 orally, in methylcobalamin form, does not seem to be effective for treating delayed sleep phase syndrome. Supplemental methylcobalamin, with or without bright light therapy, does not seem to help people with primary circadian rhythm sleep disorders.
Preliminary evidence suggests that there is no relationship between vitamin B12 status and lung cancer.
In people with a history of stroke, neither high dose vitamin B12 combinations containing pyridoxine, vitamin B12, and folic acid nor low dose combinations containing pyridoxine, vitamin B12, and folic acid seem to affect risk of recurring stroke.
Vitamin B12 is contraindicated in early Leber's disease, which is a hereditary optic nerve atrophy.
TraditionWARNING: DISCLAIMER: The below uses are based on tradition, scientific theories, or limited research. They often have not been thoroughly tested in humans, and safety and effectiveness have not always been proven. Some of these conditions are potentially serious, and should be evaluated by a qualified healthcare provider. There may be other proposed uses that are not listed below.
Aging, AIDS, allergies, amyotrophic lateral sclerosis, asthma, autism, chronic fatigue syndrome, cognitive function, depression, depressive disorder (major), diabetes, diabetic peripheral neuropathy, energy level enhancement, growth disorders (failure to thrive), hemorrhage, immunosuppression, improving concentration, inflammatory bowel disease, kidney disease, liver disease, male infertility, malignant tumors, memory loss, mood (elevate), mouth and throat inflammation (atrophic glossitis), multiple sclerosis, myoclonic disorders (spinal myoclonus), neural tube defects, osteoporosis, periodontal disease, poisoning (cyanide), protection from tobacco smoke, psychiatric disorders, seborrheic dermatitis, seizure disorders (West syndrome), tendonitis, thrombosis, thyrotoxicosis / thyroid storm (adjunct iodides), tinnitus, tremor, vitiligo.
Adults (over 18 years old)
Recommended dietary allowances (RDAs) are 2.4 micrograms per day for adults and adolescents aged 14 years and older, 2.6 micrograms per day for adult and adolescent pregnant females, and 2.8 micrograms per day for adult and adolescent lactating females. Because 10-30% of older people do not absorb food-bound vitamin B12 efficiently, those over 50 years of age should meet the RDA by eating foods fortified with B12 or by taking a vitamin B12 supplement. Supplementation of 25-100 micrograms per day has been used to maintain vitamin B12 levels in older people. A doctor and pharmacist should be consulted for use in other indications. Vitamin B12 has been taken by mouth and given by intramuscular (IM) injection by healthcare professionals. One clinical trial tested patients' acceptance of intranasal vitamin B12 replacement therapy (500 micrograms per week).
Children (under 18 years old)
Recommended dietary allowances (RDAs) have not been established for all pediatric age groups; therefore Adequate Intake (AI) levels have been used instead. The RDA and AI of vitamin B12 are: infants 0-6 months, 0.4 micrograms (AI); infants 7-12 months, 0.5 micrograms (AI); children 1-3 years, 0.9 micrograms; children 4-8 years, 1.2 micrograms; and children 9-13 years, 1.8 micrograms.
SafetyDISCLAIMER: Many complementary techniques are practiced by healthcare professionals with formal training, in accordance with the standards of national organizations. However, this is not universally the case, and adverse effects are possible. Due to limited research, in some cases only limited safety information is available.
Vitamin B12 supplements should be avoided in people sensitive or allergic to cobalamin, cobalt, or any other product ingredients.
Side Effects and Warnings
Caution should be used in patients undergoing angioplasty since an intravenous loading dose of folic acid, vitamin B6, and vitamin B12 followed by oral administration taken daily after coronary stenting might actually increase restenosis rates. Due to the potential for harm, this combination of vitamins should not be recommended for patients receiving coronary stents.
Itching, rash, transitory exanthema, and urticaria have been reported. Vitamin B12 and pyridoxine has been associated with cases of rosacea fulminans, characterized by intense erythema with nodules, papules, and pustules. Symptoms may persist for up to four months after the supplement is stopped, and may require treatment with systemic corticosteroids and topical therapy.
Diarrhea has been reported.
Peripheral vascular thrombosis has been reported. Treatment of vitamin B12 deficiency can unmask polycythemia vera, which is characterized by an increase in blood volume and the number of red blood cells. The correction of megaloblastic anemia with vitamin B12 can result in fatal hypokalemia and gout in susceptible individuals, and it can obscure folate deficiency in megaloblastic anemia. Caution is warranted.
Pregnancy and Breastfeeding
Vitamin B12 is likely safe when used orally in amounts that do not exceed the recommended dietary allowance (RDA).
There is insufficient reliable information available about the safety of larger amounts of vitamin B12 during pregnancy.
Interactions with Drugs
Excessive alcohol intake lasting longer than two weeks can decrease vitamin B12 absorption from the gastrointestinal tract.
Aminosalicylic acid can reduce oral vitamin B12 absorption, possibly by as much as 55%, as part of a general malabsorption syndrome. Megaloblastic changes, and occasional cases of symptomatic anemia, have occurred. Vitamin B12 levels should be monitored in people taking aminosalicylic acid for more than one month.
An increased bacterial load can bind significant amounts of vitamin B12 in the gut, preventing its absorption. In people with bacterial overgrowth of the small bowel, antibiotics such as metronidazole (Flagyl®) can actually improve vitamin B12 status. The effects of most antibiotics on gastrointestinal bacteria are unlikely to have clinically significant effects on vitamin B12 levels.
The data regarding the effects of oral contraceptives on vitamin B12 serum levels are conflicting. Some studies have found reduced serum levels in birth control pill users, but others have found no effect despite the use of birth control pills for up to six months. When birth control pill use is stopped, normalization of vitamin B12 levels usually occurs. Lower vitamin B12 serum levels seen with birth control pills probably are not clinically significant.
Limited case reports suggest that chloramphenicol can delay or interrupt the reticulocyte response to supplemental vitamin B12 in some patients. Blood counts should be monitored closely if this combination cannot be avoided.
Cobalt irradiation of the small bowel can decrease gastrointestinal (GI) absorption of vitamin B12.
Colchicine can disrupt normal intestinal mucosal function, leading to malabsorption of several nutrients, including vitamin B12. Lower doses do not seem to have a significant effect on vitamin B12 absorption after three years of colchicine therapy. The significance of this interaction is unclear. Vitamin B12 levels should be monitored in people taking large doses of colchicine for prolonged periods.
Colestipol (Colestid®) and Cholestyramine (Questran®) resins can decrease gastrointestinal (GI) absorption of vitamin B12. It is unlikely that this interaction will deplete body stores of vitamin B12 unless there are other factors contributing to deficiency. In a group of children treated with cholestyramine for up to 2.5 years, there was not any change in serum vitamin B12 levels. Routine supplements are not necessary.
H2-blockers include cimetidine (Tagamet®), famotidine (Pepcid®), nizatidine (Axid®), and ranitidine (Zantac®). Reduced secretion of gastric acid and pepsin produced by H2-blockers can reduce absorption of protein-bound (dietary) vitamin B12, but not of supplemental vitamin B12. Gastric acid is needed to release vitamin B12 from protein for absorption. Clinically significant vitamin B12 deficiency and megaloblastic anemia are unlikely, unless H2-blocker therapy is prolonged (two years or more) or the person's diet is poor. It is also more likely if the person is rendered achlorhydric (lacking hydrochloric stomach acid), which occurs more frequently with proton pump inhibitors than H2-blockers. Vitamin B12 levels should be monitored in people taking high doses of H2 blockers for prolonged periods.
Metformin may reduce serum folic acid and vitamin B12 levels. These changes can lead to hyperhomocysteinemia (abnormally large levels of homocysteine in the blood), adding to the risk of cardiovascular disease in people with diabetes. There are also rare reports of megaloblastic anemia in people who have taken metformin for five years or more. Reduced serum levels of vitamin B12 occur in up to 30% of people taking metformin chronically. However, clinically significant deficiency is not likely to develop if dietary intake of vitamin B12 is adequate. Deficiency can be corrected with vitamin B12 supplements even if metformin is continued. The metformin-induced malabsorption of vitamin B12 is reversible by oral calcium supplementation. A multivitamin preparation may also be valuable for some patients. Patients should be monitored for signs and symptoms of vitamin B12 and folic acid deficiency. People taking metformin chronically should be advised to include adequate amounts of vitamin B12 in their diet, and have their serum vitamin B12 and homocysteine levels checked annually.
Nicotine can reduce serum vitamin B12 levels. The need for vitamin B12 supplementation has not been adequately studied.
Nitrous oxide inactivates the cobalamin form of vitamin B12 by oxidation. Symptoms of vitamin B12 deficiency, including sensory neuropathy, myelopathy, and encephalopathy can occur within days or weeks of exposure to nitrous oxide anesthesia in people with subclinical vitamin B12 deficiency. Symptoms are treated with high doses of vitamin B12, but recovery can be slow and incomplete. People with normal vitamin B12 levels have sufficient vitamin B12 stores to make the effects of nitrous oxide insignificant, unless exposure is repeated and prolonged (nitrous oxide abuse). Vitamin B12 levels should be checked in people with risk factors for vitamin B12 deficiency prior to using nitrous oxide anesthesia.
Phenytoin (Dilantin®), phenobarbital, and primidone (Mysoline®) anticonvulsants have been associated with reduced vitamin B12 absorption and reduced serum and cerebrospinal fluid levels in some patients. This may contribute to the megaloblastic anemia, primarily caused by folate deficiency, associated with these drugs. It has also been suggested that reduced vitamin B12 levels may contribute to the neuropsychiatric side effects of these drugs. Patients should be encouraged to maintain adequate dietary vitamin B12 intake. Folate and vitamin B12 status should be checked if symptoms of anemia develop.
Proton pump inhibitors (PPIs) include omeprazole (Prilosec®, Losec®), lansoprazole (Prevacid®), rabeprazole (Aciphex®), pantoprazole (Protonix®, Pantoloc®), and esomeprazole (Nexium®). The reduced secretion of gastric acid and pepsin produced by PPIs can reduce absorption of protein-bound (dietary) vitamin B12, but not supplemental vitamin B12. Gastric acid is needed to release vitamin B12 from protein for absorption. Reduced vitamin B12 levels may be more common with PPIs than with H2-blockers, because they are more likely to produce achlorhydria (complete absence of gastric acid secretion). However, clinically significant vitamin B12 deficiency is unlikely, unless PPI therapy is prolonged (two years or more) or dietary vitamin intake is low. Vitamin B12 levels should be monitored in people taking high doses of PPIs for prolonged periods.
Reduced serum vitamin B12 levels may occur when zidovudine (AZT, Combivir®, Retrovir®) therapy is started. This adds to other factors that cause low vitamin B12 levels in people with HIV and might contribute to the hematological toxicity associated with zidovudine. However, data suggests vitamin B12 supplements are not helpful for people taking zidovudine.
Interactions with Herbs and Dietary Supplements
Folic acid, particularly in large doses, can mask vitamin B12 deficiency. In vitamin B12 deficiency, folic acid can produce hematologic improvement in megaloblastic anemia, while allowing potentially irreversible neurological damage to progress. Vitamin B12 status should be determined before folic acid is given as a monotherapy.
Potassium supplements can reduce absorption of vitamin B12 in some people. This effect has been reported with potassium chloride and, to a lesser extent, with potassium citrate. Potassium might contribute to vitamin B12 deficiency in some people with other risk factors, but routine supplements are not necessary.
Preliminary evidence suggests that vitamin C supplements can destroy dietary vitamin B12. However, other components of food, such as iron and nitrates, might counteract this effect. Clinical significance is unknown, and it can likely be avoided if vitamin C supplements are taken at least two hours after meals.
This information is based on a systematic review of scientific literature and consensus statements edited and peer-reviewed by contributors to the Natural Standard Research Collaboration (www.naturalstandard.com): Ethan Basch, MD (Memorial Sloan-Kettering Cancer Center); Dawn Costa, BA, BS (Natural Standard Research Collaboration); Jill M. Grimes Serrano, PhD (Natural Standard Research Collaboration); Jenna Hollenstein, MS, RD (Natural Standard Research Collaboration); Shaina Tanguay-Colucci, BS (Natural Standard Research Collaboration); Catherine Ulbricht, PharmD (Massachusetts College of Pharmacy); Christine Ulbricht, PharmD (University of Massachusetts); Wendy Weissner, BA (Natural Standard Research Collaboration).
BibliographyDISCLAIMER: Natural Standard developed the above evidence-based information based on a thorough systematic review of the available scientific articles. For comprehensive information about alternative and complementary therapies on the professional level, go to www.naturalstandard.com. Selected references are listed below.
Albert CM, Cook NR, Gaziano JM, et al. Effect of folic acid and B vitamins on risk of cardiovascular events and total mortality among women at high risk for cardiovascular disease: a randomized trial. JAMA. 2008 May 7;299(17):2027-36.
Andres E, Kurtz JE, Perrin AE, et al. Oral cobalamin therapy for the treatment of patients with food-cobalamin malabsorption. Am J Med 2001;111:126-9.
Bjørke-Monsen AL, Torsvik I, Saetran H, et al. Common metabolic profile in infants indicating impaired cobalamin status responds to cobalamin supplementation. Pediatrics 2008 Jul;122(1):83-91.
Erol I, Alehan F, Gümüs A. West syndrome in an infant with vitamin B12 deficiency in the absence of macrocytic anaemia. Dev Med Child Neurol 2007 Oct;49(10):774-6.
Eussen SJ, de Groot LC, Clarke R, et al. Oral cyanocobalamin supplementation in older people with vitamin B12 deficiency: a dose-finding trial. Arch Intern Med 2005 May 23;165(10):1167-72.
Force RW, Meeker AD, Cady PS, et al. Increased vitamin B12 requirement associated with chronic acid suppression therapy. Ann Pharmacother 2003;37:490-3.
Haggarty P, McCallum H, McBain H, Effect of B vitamins and genetics on success of in-vitro fertilisation: prospective cohort study. Lancet 2006 May 6;367(9521):1513-9.
Lehman JS, Bruce AJ, Rogers RS. Atrophic glossitis from vitamin B12 deficiency: a case misdiagnosed as burning mouth disorder. J Periodontol 2006 Dec;77(12):2090-2.
Malouf R, Grimley Evans J. Folic acid with or without vitamin B12 for the prevention and treatment of healthy elderly and demented people. Cochrane Database Syst Rev 2008 Oct 8;(4):CD004514.
Molloy AM, Kirke PN, Brody LC, et al. Effects of folate and vitamin B12 deficiencies during pregnancy on fetal, infant, and child development. Food Nutr Bull 2008 Jun;29(2 Suppl):S101-11; discussion S112-5.
Ryan-Harshman M, Aldoori W. Vitamin B12 and health. Can Fam Physician 2008 Apr;54(4):536-41.
Seal EC, Metz J, Flicker L, et al. A randomized, double-blind, placebo-controlled study of oral vitamin B12 supplementation in older patients with subnormal or borderline serum vitamin B12 concentrations. J Am Geriatr Soc 2002;50:146-51.
Suzuki DM, Alagiakrishnan K, Masaki KH, et al. Patient acceptance of intranasal cobalamin gel for vitamin B12 replacement therapy. Hawaii Med J 2006 Nov;65(11):311-4.
Toole JF, Malinow MR, Chambless LE, et al. Lowering homocysteine in patients with ischemic stroke to prevent recurrent stroke, myocardial infarction, and death: the Vitamin Intervention for Stroke Prevention (VISP) randomized controlled trial. JAMA 2004;291:565-75.
Yajnik CS, Lubree HG, Thuse NV, et al. Oral vitamin B12 supplementation reduces plasma total homocysteine concentration in women in India. Asia Pac J Clin Nutr 2007;16(1):103-9.
Remember, keep this and all other medicines out of the reach of children, never share your medicines with others, and use this medication only for the indication prescribed. |
McDonnell Douglas F-15 Eagle
|USAF F-15C during an Operation Noble Eagle patrol|
|Role||Air superiority fighter|
Boeing Defense, Space & Security
|First flight||27 July 1972|
|Introduction||9 January 1976|
|Primary users||United States Air Force
Japan Air Self-Defense Force
Royal Saudi Air Force
Israeli Air Force
|Number built||F-15A/B/C/D/J/DJ: 1,198|
|Unit cost||F-15A/B: US$28 million (1998)
F-15C/D: US$30 million (1998)
|Variants||McDonnell Douglas F-15E Strike Eagle
McDonnell Douglas F-15 STOL/MTD
Boeing F-15SE Silent Eagle
The McDonnell Douglas (now Boeing) F-15 Eagle is a twin-engine, all-weather tactical fighter designed by McDonnell Douglas to gain and maintain air superiority in aerial combat. It is considered among the most successful modern fighters, with over 100 aerial combat victories with no losses in dogfights. Following reviews of proposals, the United States Air Force selected McDonnell Douglas' design in 1967 to meet the service's need for a dedicated air superiority fighter. The Eagle first flew in July 1972, and entered service in 1976.
Since the 1970s, the Eagle has been exported to Israel, Japan, Saudi Arabia, and other nations. The F-15 was originally envisioned as a pure air superiority aircraft. Its design included a secondary ground-attack capability that was largely unused. The design proved flexible enough that an all-weather strike derivative, the F-15E Strike Eagle, was later developed, and entered service in 1989. The F-15 Eagle is expected to be in service with the U.S. Air Force past 2025.
Following studies in 1964–1965, the U.S. Air Force developed requirements for an air superiority fighter in October 1965. Then on 8 December 1965, the service issued a request for proposals (RFP) for the new fighter. The request called for both air-to-air and air-to-ground capabilities. Eight companies responded with proposals. In the following study phase, four of these companies developed some 500 design concepts. Typical designs featured variable-sweep wings, weighed over 60,000 lb (27,200 kg), included a top speed of Mach 2.7 and a thrust-to-weight ratio of 0.75. The designs were not accepted by the Air Force as they compromised fighter qualities for ground attack qualities. Acceptance of the Energy-Maneuverability (E-M) theory by the Air Force led to a change in requirements for improved maneuverability by the spring 1967. The design mission weight was reduced to 40,000 lb (18,100 kg), top speed reduced to Mach 2.3–2.5 and thrust-to-weight ratio increased to 0.97.
In 1967 U.S. intelligence was surprised to find that the Soviet Union was producing a large fighter aircraft, the MiG-25 'Foxbat'. It was not known in the West at the time that the MiG-25 was designed as a high-speed interceptor, not an air superiority fighter, so its primary asset was speed, not maneuverability. The MiG-25's huge tailplanes and vertical stabilizers (tail fins) hinted at a very maneuverable aircraft, which worried the Air Force that its performance might be better than its U.S. counterparts. In reality, the MiG's large fins and stabilators were necessary to prevent the aircraft from encountering inertia coupling in high-speed, high-altitude flight.
The McDonnell Douglas F-4 Phantom II of the USAF, U.S. Navy and U.S. Marine Corps was the only fighter with enough power, range, and maneuverability to be given the primary task of dealing with the threat of Soviet fighters while flying with visual engagement rules. As a matter of policy, the Phantoms could not engage targets without positive visual identification, so they could not engage targets at long ranges, as designed. Medium-range AIM-7 Sparrow missiles, and to a lesser degree even the AIM-9 Sidewinder, were often unreliable and ineffective at close ranges where it was found that guns were often the only effective weapon. The Phantom did not originally have any guns or cannons, but experience in Vietnam led to the addition of an internally mounted cannon in later versions.
F-X program
There was a clear need for a new fighter that overcame the close-range limitation of the Phantom while retaining long-range air superiority. After rejecting the U.S. Navy VFX program (which led to the F-14 Tomcat) as being unsuited to its needs, the U.S. Air Force issued its own requirements for the F-X (read as Fighter-Unknown, sometimes referred to as Fighter-Experimental), a specification for a relatively lightweight air superiority fighter. The requirements called for single-seat fighter having a maximum take-off weight of 40,000 lb (18,100 kg) for the air-to-air role with a maximum speed of Mach 2.5 and a thrust to weight ratio of nearly 1 at mission weight. Four companies submitted proposals, with the Air Force eliminating General Dynamics and awarding contracts to Fairchild Republic, North American Rockwell, and McDonnell Douglas for the definition phase in December 1968. The companies submitted technical proposals by June 1969. The Air Force announced the selection of McDonnell Douglas on 23 December 1969. The winning design resembled the twin-tailed F-14, but with fixed wings. It would not be significantly lighter or smaller than the F-4 that it would replace.
The Eagle's initial versions were the F-15 single-seat variant and TF-15 twin-seat variant. (After the F-15C was first flown in 1980 the designations were changed to "F-15A" and "F-15B"). These versions would be powered by new Pratt & Whitney F100 engines to achieve a combat thrust-to-weight ratio in excess of 1. A proposed 25 mm Ford-Philco GAU-7 cannon with caseless ammunition suffered development problems. It was dropped in favor of the standard M61 Vulcan gun. The F-15 used conformal carriage of four Sparrow missiles like the Phantom. The fixed wing was put onto a flat, wide fuselage that also provided an effective lifting surface. The first F-15A flight was made in July 1972 with the first flight of the two-seat F-15B following in July 1973.
The F-15 has a "look-down/shoot-down" radar that can distinguish low-flying moving targets from ground clutter. The F-15 would use computer technology with new controls and displays to lower pilot workload and require only one pilot to save weight. Unlike the F-14 or F-4, the F-15 has only a single canopy frame with clear vision forward. The USAF introduced the F-15 as "the first dedicated USAF air superiority fighter since the North American F-86 Sabre."
The F-15 was favored by customers such as the Israel and Japan air arms. Criticism from the fighter mafia that the F-15 was too large to be a dedicated dogfighter, and too expensive to procure in large numbers, led to the Lightweight Fighter (LWF) program, which led to the USAF General Dynamics F-16 Fighting Falcon and the middle-weight Navy McDonnell Douglas F/A-18 Hornet.
Further development
The single-seat F-15C and two-seat F-15D models entered production in 1978 and conducted their first flights in February and June of that year. These models were fitted with the Production Eagle Package (PEP 2000), which included 2,000 lb (900 kg) of additional internal fuel, provisions for carrying exterior conformal fuel tanks, and an increased maximum takeoff weight of up to 68,000 lb (30,700 kg). The increased takeoff weight allows internal fuel, a full weapons load, conformal fuel tanks, and three external fuel tanks to be carried. The APG-63 radar uses a programmable signal processor (PSP), enabling the radar to be reprogrammable for additional purposes such as the addition of new armaments and equipment. The PSP was the first of its kind in the world, and the upgraded APG-63 radar was the first radar to use it. Other improvements on the C and D models included strengthened landing gear, a new digital central computer, and an overload warning system, which allows the pilot to fly the fighter to 9 g at all weights.
The F-15 Multistage Improvement Program (MSIP) was initiated in February 1983 with the first production MSIP F-15C produced in 1985. Improvements included an upgraded central computer; a Programmable Armament Control Set, allowing for advanced versions of the AIM-7, AIM-9, and AIM-120A missiles; and an expanded Tactical Electronic Warfare System that provides improvements to the ALR-56C radar warning receiver and ALQ-135 countermeasure set. The final 43 F-15Cs included the enhanced-capability Hughes APG-70 radar, which was developed for the F-15E. These 43 F-15Cs with APG-70 radar are sometimes referred as Enhanced Eagles. Earlier MSIP F-15Cs with the APG-63 were upgraded to the APG-63(V)1 to significantly improve maintainability and enable performance similar to the APG-70. Existing F-15s were retrofitted with these improvements.
In 1979, McDonnell Douglas and F-15 radar manufacturer, Hughes, teamed to privately develop a strike fighter version of the F-15. This version competed in the Air Force's Dual-Role Fighter competition starting in 1982. The Boeing F-15E strike variant was selected for production in 1984. Beginning in 1985, F-15C and D models were equipped with the improved P&W F100-220 engine and digital engine controls, providing quicker throttle response, reduced wear, and lower fuel consumption. Starting in 1997, original F100-100 engines were upgraded to a similar configuration with the designation F100-220E starting.
Beginning in 2007, 178 USAF F-15Cs were retrofitted with the AN/APG-63(V)3 Active Electronically Scanned Array (AESA) radar. A significant number of F-15s are to be equipped with the Joint Helmet Mounted Cueing System (JHMCS). Lockheed Martin is working on an IRST system for the F-15C. A follow-on upgrade called the Eagle passive/active warning survivability system (EPAWSS) was planned but remained unfunded.
The F-15 has an all-metal semi-monocoque fuselage with a large cantilever shoulder-mounted wing. The empennage is metal and composite construction, with twin aluminum/composite honeycomb fins with boron-composite skins, resulting in an exceptionally thin tailplane and rudders with all-moving composite horizontal tail surfaces outboard of the fins. The F-15 has a spine-mounted air brake and retractable tricycle landing gear. It is powered by two Pratt & Whitney F100 axial-flow turbofan engines with afterburners mounted side-by-side in the fuselage. The cockpit is mounted high in the forward fuselage with a one-piece windscreen and large canopy to increase visibility.
The F-15's maneuverability is derived from low wing loading (weight to wing area ratio) with a high thrust-to-weight ratio enabling the aircraft to turn tightly without losing airspeed. The F-15 can climb to 30,000 ft (10,000 m) in around 60 seconds. The thrust output of the dual engines is greater than the aircraft's weight, thus giving it the ability to accelerate in a vertical climb. The weapons and flight control systems are designed so that one person can safely and effectively perform air-to-air combat. The A and C-models are single-seat variants that make up the bulk of F-15 production. B and D-models add a second seat behind the pilot for training. E-models use the second seat for a bombardier/navigator. Visibly, the F-15 has a unique feature vis a vis other modern fighter aircraft in that it does not have the distinctive turkey feather aerodynamic exhaust petals covering its engine nozzles. This is because the petal design on the F-15 was problematic and could fall off in flight; therefore they were removed, resulting in a 3% drag increase.
A multi-mission avionics system includes a heads-up display (HUD), advanced radar, inertial guidance system (INS), flight instruments, ultra high frequency (UHF) communications, and Tactical Air Navigation (TACAN) and Instrument Landing System (ILS) receivers. It also has an internally mounted, tactical electronic-warfare system, identification, friend or foe (IFF) system, electronic countermeasures suite and a central digital computer.
The heads-up display projects, through a combiner, all essential flight information gathered by the integrated avionics system. This display, visible in any light condition, provides the pilot information necessary to track and destroy an enemy aircraft without having to look down at cockpit instruments.
The F-15's versatile APG-63/70 pulse-Doppler radar system can look up at high-flying targets and down at low-flying targets without being confused by ground clutter. It can detect and track aircraft and small high-speed targets at distances beyond visual range (the maximum being 120 nautical miles (220 km) away) down to close range, and at altitudes down to treetop level. The radar feeds target information into the central computer for effective weapons delivery. The capability of locking onto targets as far as 50 nautical miles (90 km) with an AIM-120 AMRAAM missile enables true beyond visual range (BVR) engagement of targets. For close-in dogfights, the radar automatically acquires enemy aircraft, and this information is projected on the heads-up display. The F-15's electronic warfare system provides both threat warning and automatic countermeasures against selected threats.
Weaponry and external stores
A variety of air-to-air weaponry can be carried by the F-15. An automated weapon system enables the pilot to perform aerial combat safely and effectively, using the head-up display and the avionics and weapons controls located on the engine throttles or control stick. When the pilot changes from one weapon system to another, visual guidance for the required weapon automatically appears on the head-up display.
The Eagle can be armed with combinations of four different air-to-air weapons: AIM-7F/M Sparrow missiles or AIM-120 AMRAAM advanced medium range air-to-air missiles on its lower fuselage corners, AIM-9L/M Sidewinder or AIM-120 AMRAAM missiles on two pylons under the wings, and an internal M61A1 20 mm Gatling gun in the right wing root.
Low-drag conformal fuel tanks (CFTs) were developed for the F-15C and D models. They can be attached to the sides of the engine air intake trunks under each wing and are designed to the same load factors and airspeed limits as the basic aircraft. They degrade performance by increasing drag and cannot be jettisoned in-flight (unlike conventional external tanks). Each conformal fuel tank can hold 750 U.S. gallons (2,840 L) of fuel. These tanks increase range and reduce the need for in-flight refueling. All external stations for munitions remain available with the tanks in use. Moreover, Sparrow or AMRAAM missiles can be attached to the corners of the conformal fuel tanks. The 57 FIS based at Keflavik NAS, Iceland was the only C-model squadron to use CFTs on a regular basis due to its extended operations over the North Atlantic. With the closure of the 57 FIS, the F-15E is the only variant to carry them on a routine basis. CFTs have also been sold to Israel and Saudi Arabia.
The F-15E Strike Eagle is a two-seat, dual-role, totally integrated fighter for all-weather, air-to-air and deep interdiction missions. The rear cockpit is upgraded to include four multi-purpose CRT displays for aircraft systems and weapons management. The digital, triple-redundant Lear Siegler flight control system permits coupled automatic terrain following, enhanced by a ring-laser gyro inertial navigation system. For low-altitude, high-speed penetration and precision attack on tactical targets at night or in adverse weather, the F-15E carries a high-resolution APG-70 radar and LANTIRN pods to provide thermal imagery.
The APG-63(V)2 Active Electronically Scanned Array (AESA) radar has been retrofitted to 18 U.S. Air Force F-15C aircraft. This upgrade includes most of the new hardware from the APG-63(V)1, but adds an AESA to provide increased pilot situational awareness. The AESA radar has an exceptionally agile beam, providing nearly instantaneous track updates and enhanced multi-target tracking capability. The APG-63(V)2 is compatible with current F-15C weapon loads and enables pilots to take full advantage of AIM-120 capabilities, simultaneously guiding multiple missiles to several targets widely spaced in azimuth, elevation, or range.
Operational history
Introduction and early service
The largest operator of the F-15 is the United States Air Force. The first Eagle (F-15B) was delivered 13 November 1974. In January 1976, the first Eagle destined for a combat squadron, the 555th TFS, was delivered. These initial aircraft carried the Hughes Aircraft (now Raytheon) APG-63 radar.
The first F-15 kill was scored by Israeli Air Force (IAF) ace Moshe Melnik in 1979. In 1979–81, during Israeli raids against Palestinian factions based in Lebanon, F-15As downed 13 Syrian MiG-21 "Fishbeds" and two Syrian MiG-25 "Foxbats", the latter being the aircraft the F-15 was designed to kill. Israeli F-15As and Bs participated as escorts in Operation Opera and served during the 1982 Lebanon War. During the latter, Israeli F-15s shot down 40 Syrian jet fighters (23 MiG-21 "Fishbeds" and 17 MiG-23 "Floggers") and one Syrian SA.342L Gazelle helicopter. Later during 1985, IAF Eagles, in Operation Wooden Leg, bombed the PLO headquarters in Tunisia. This was one of the few times air superiority F-15s (A/B/C/D models) were used in tactical strike missions. The air to ground role for the air superiority variants became more frequently used in Israeli service starting from the opening years of the new century with the integration of GPS guided bombs.
Anti-satellite trials
The ASM-135 missile was designed to be a standoff anti-satellite (ASAT) weapon, with the F-15 acting as a first stage. The Soviet Union could correlate a U.S. rocket launch with a spy satellite loss, but an F-15 carrying an ASAT would blend in among hundreds of F-15 flights. From January 1984 to September 1986, two F-15As were used as launch platforms for the ASAT missile. The F-15As were modified to carry one ASM-135 on the centerline station with extra equipment within a special centerline pylon. The launch aircraft executed a Mach 1.22, 3.8 g climb at 65° to release the ASAT missile at an altitude of 38,100 ft (11,600 m). The flight computer was updated to control the zoom-climb and missile release.
The third test flight involved a retired P78-1 solar observatory satellite in a 345-mile (555 km) orbit, which was destroyed by kinetic energy. The pilot, USAF Major Wilbert D. "Doug" Pearson, became the only pilot to destroy a satellite. The ASAT program involved five test launches. The program was officially terminated in 1988.
Gulf War and aftermath
The USAF began deploying F-15C, D and E model aircraft to the Persian Gulf region in August 1990 for Operations Desert Shield and Desert Storm. During Gulf War combat against Iraqi forces, they accounted for 36 of the 39 Air Force air-to-air victories. F-15Es were operated mainly at night, hunting modified SCUD missile launchers and artillery sites using the LANTIRN system. According to the USAF, its F-15Cs had 34 confirmed kills of Iraqi aircraft during the 1991 Gulf War, mostly by missile fire: five MiG-29 "Fulcrums", two MiG-25 "Foxbats", eight MiG-23 "Floggers", two MiG-21 "Fishbeds", two Su-25 "Frogfoots", four Su-22 "Fitters", one Su-7, six Mirage F1s, one Il-76 cargo plane, one Pilatus PC-9 trainer, and two Mi-8 helicopters. Air superiority was achieved in the first three days of the conflict; many of the later kills were reportedly of Iraqi aircraft fleeing to Iran, rather than trying to engage U.S. aircraft. The single-seat F-15C was used for air superiority, and the F-15E was heavily used in air-to-ground attacks. An F-15E achieved an aerial kill of another Iraqi Mi-8 helicopter using a laser-guided bomb during the air war. The F-15E sustained two losses to ground fire in the Gulf War in 1991. Another one was damaged on the ground by a SCUD strike on Dhahran air base.
They have since been deployed to support Operation Southern Watch, the patrolling of the No-Fly Zone in Southern Iraq; Operation Provide Comfort in Turkey; in support of NATO operations in Bosnia, and recent air expeditionary force deployments. In 1994, two U.S. Army Sikorsky UH-60 Black Hawks were mistakenly downed by USAF F-15Cs in northern Iraq in a friendly fire incident. USAF F-15Cs shot down four Yugoslav MiG-29s using AIM-120 missiles during NATO's 1999 intervention in Kosovo, Operation Allied Force.
Structural defects
All F-15 aircraft were grounded by the US Air Force after a Missouri Air National Guard F-15C came apart in flight and crashed on 2 November 2007. The newer F-15E fleet was later cleared for continued operations. The US Air Force reported on 28 November 2007 that a critical location in the upper longerons on the F-15C model was suspected of causing the failure, causing the fuselage forward of the air intakes, including the cockpit and radome, to separate from the airframe.
F-15A through D-model aircraft were ordered grounded until the location received more detailed inspections and repairs as needed. The grounding of F-15s received media attention as it began to place strains on the nation's air defense efforts. The grounding forced some states to rely on their neighboring states' fighters for air defense protection, and Alaska to depend on Canadian Forces' fighter support.
On 8 January 2008, the USAF Air Combat Command (ACC) cleared a portion of its F-15A through D-model fleet for return to flying status. It also recommended a limited return to flight for units worldwide using the affected models. The accident review board report was released on 10 January 2008. The report stated that analysis of the F-15C wreckage determined that the longeron did not meet drawing specifications, which led to fatigue cracks and finally a catastrophic failure of the remaining support structures and breakup of the aircraft in flight. In a report released on 10 January 2008, nine other F-15s were identified to have similar problems in the longeron. As a result of these problems, General John D. W. Corley stated that "the long-term future of the F-15 is in question." On 15 February 2008, ACC cleared all its grounded F-15A/B/C/D fighters for flight pending inspections, engineering reviews and any needed repairs. ACC also recommended release of other U.S. F-15A/B/C/D aircraft.
Recent service
Indian Air Force (IAF) Su-30MKs, MiG-29s and other fighters achieved success in air combat exercises against U.S. Air Force (USAF) F-15Cs during "Cope India" in February 2004. The U.S. agreed not to use beyond visual range AIM-120 missiles on its fighters. During the USAF's Red Flag advanced combat training exercises in 2008, American F-15Cs, F-16Cs, and F-22s bested Indian Su-30MKIs. The Su-30s operated with their radars in training mode to keep their signals secret.
The F-15 in all air forces had a combined air-to-air combat record of 104 kills to 0 losses as of February 2008. No air superiority versions of the F-15 (A/B/C/D models) have been shot down by enemy forces. Over half of F-15 kills were achieved by Israeli Air Force pilots.
With the retirement of the F-15A and B models, the F-15C and D models are being supplemented in U.S. service by the F-22 Raptor. As of 2013, Regular Air Force F-15C and F-15D fighters are based overseas with the Pacific Air Forces (PACAF) at Kadena AB in Japan and with the U.S. Air Forces in Europe (USAFE) at RAF Lakenheath in the United Kingdom. Other Regular Air Force F-15s are operated by Air Combat Command as adversary/aggressor platforms at Nellis AFB, Nevada, and by Air Force Material Command in test and evaluation roles at Edwards AFB, California and Eglin AFB, Florida. All remaining combat coded F-15Cs and F-15Ds are operated by the Air National Guard.
USAF is upgrading 178 F-15C/Ds with the AN/APG-63(V)3 AESA radar, and upgrade other F-15s with the Joint Helmet Mounted Cueing System. The Air Force plans to keep 178 F-15C/Ds along with 224 F-15Es in service beyond 2025. The F-15E will remain in service for years to come because of the model's primary air-to-ground role and the lower number of hours on the F-15E airframes.
Basic models
- Single-seat all-weather air-superiority fighter version, 384 built 1972–1979.
- Two-seat training version, formerly designated TF-15A, 61 built 1972–1979.
- Improved single-seat all-weather air-superiority fighter version, 483 built 1979–1985. The last 43 F-15Cs are being upgraded with AN/APG-70 radar.
- Two-seat training version, 92 built 1979–1985.
- Single-seat all-weather air-superiority fighter version for the Japan Air Self-Defense Force 139 built under license in Japan by Mitsubishi 1981–1997, two built in St. Louis.
- Two-seat training version for the Japan Air Self-Defense Force. 12 built in St. Louis, and 25 built under license in Japan by Mitsubishi during 1981–1997.
- F-15N Sea Eagle
- The F-15N was a carrier-capable variant proposed in the early 1970s to the U.S. Navy as an alternative to the heavier and, at the time, considered as "riskier" technology program: Grumman F-14 Tomcat. The F-15N-PHX was another proposed naval version capable of carrying the AIM-54 Phoenix missile. These featured folding wingtips, reinforced landing gear and a stronger tailhook for shipboard operation.
- F-15E Strike Eagle
- See McDonnell Douglas F-15E Strike Eagle for F-15E, F-15I, F-15S, F-15K, F-15SG, and other F-15E-based variants.
- F-15SE Silent Eagle
- See Boeing F-15SE Silent Eagle for a recent proposed F-15E variant with a reduced radar signature.
Twelve prototypes were built and were used for trials by the F-15 Joint Test Force at Edwards Air Force Base using McDonnell Douglas and United States Air Force personnel. Many of the prototypes were later used by NASA for trials and experiments.
- F-15A-1, AF Ser. No. 71-0280
- Was the first F-15 to fly on 11 July 1972 from Edwards Air Force Base, it was used as a trial aircraft for exploring the flight envelope, general handling and testing the carriage of external stores.
- F-15A-1, AF Ser. No. 71-0281
- The second prototype first flew on 26 September 1972 and was used to test the F100 engine.
- F-15A-2, AF Ser. No. 71-0282
- First flew on 4 November 1972 and was used to test the APG-62 radar and avionics.
- F-15A-2, AF Ser. No. 71-0283
- First flew on 13 January 1973 and was used as a structural test aircraft, it was the first aircraft to have the smaller wingtips to clear a severe buffet problem found on earlier aircraft.
- F-15A-2, AF Ser. No. 71-0284
- First flew on 7 March 1973 it was used for armament development and was the first aircraft fitted with an internal cannon.
- F-15A-3, AF Ser. No. 71-0285
- First flew on 23 May 1973 and was used to test the missile fire control system and other avionics.
- F-15A-3, AF Ser. No. 71-0286
- First flew on 14 June 1973 and was used for armament trials and testing external fuel stores.
- F-15A-4, AF Ser. No. 71-0287
- First flew on 25 August 1973 and was used for spin recovery, angle of attack and fuel system testing, it was fitted with an anti-spin recovery parachute. The aircraft was loaned to NASA from 1976 for engine development trials.
- F-15A-4, AF Ser. No. 71-0288
- First flew on 20 October 1973 and was used to test integrated aircraft and engine performance, it was later used by McDonnell Douglas as a test aircraft in the 1990s.
- F-15A-4, AF Ser. No. 71-0289
- First flew on 30 January 1974 and was used for trials on the radar, avionics and electronic warfare systems.
- F-15B-3, AF Ser. No. 71-0290
- The first two-seat prototype originally designated the TF-15A, it first flew on 7 July 1973.
- F-15B-4, AF Ser. No. 71-0291
- First flew on 18 October 1973 as a TF-15A and used as a test and demonstration aircraft. In 1976 it made an overseas sales tour painted in markings to celebrate the bicentenary of the United States.
Research and test
- F-15 Streak Eagle (AF Ser. No.72-0119)
- One stripped of most avionics and unpainted F-15A, demonstrated the fighter's acceleration – broke eight time-to-climb world records between 16 January and 1 February 1975 at Grand Forks AFB, ND. It was delivered to the National Museum of the United States Air Force in December 1980.
- F-15 STOL/MTD (AF Ser. No. 71-0290)
- The first F-15B was converted into a short takeoff and landing, maneuver technology demonstrator aircraft. In the late 1980s it received canard flight surfaces in addition to its usual horizontal tail, along with square thrust-vectoring nozzles. It was used as a short-takeoff/maneuver-technology (SMTD) demonstrator.
- F-15 ACTIVE (AF Ser. No. 71-0290)
- The F-15 S/MTD was later converted into an advanced flight control technology research aircraft with thrust vectoring nozzles.
- F-15 IFCS (AF Ser. No. 71-0290)
- The F-15 ACTIVE was then converted into an intelligent flight control systems research aircraft. F-15B 71-0290 was the oldest F-15 still flying when retired in January 2009.
- F-15 MANX
- Concept name for a tailless variant of the F-15 ACTIVE, but the NASA ACTIVE experimental aircraft was never modified to be tailless.
- F-15 Flight Research Facility (AF Ser. No. 71-0281 and AF Ser. No. 71-0287)
- Two F-15A aircraft were acquired in 1976 for use by NASA's Dryden Flight Research Center for numerous experiments such as: Highly Integrated Digital Electronic Control (HiDEC), Adaptive Engine Control System (ADECS), Self-Repairing and Self-Diagnostic Flight Control System (SRFCS) and Propulsion Controlled Aircraft System (PCA). 71-0281, the second flight-test F-15A, was returned to the Air Force and became a static display at Langley AFB in 1983.
- F-15B Research Testbed (AF Ser. No. 74-0141)
- Acquired in 1993, it was an F-15B modified and used by NASA's Dryden Flight Research Center for flight tests.
- Israeli Air Force has operated F-15s since 1977. The IAF has 43 F-15A/B/C/D (20 F-15A, 6 F-15B, 11 F-15C, and 6 F-15D) aircraft in service as of January 2011.
- Saudi Arabia
- Royal Saudi Air Force has 70 F-15C/D (49 F-15C and 21 F-15D) Eagles in operation as of January 2011.
- United States Air Force operates 254 F-15C/D aircraft (114 Regular Air Force and 140 Air National Guard) as of September 2010.
Notable accidents
- On 1 May 1983, during an Israeli Air Force training dogfight, an F-15D collided with a Douglas A-4 Skyhawk. Unknown to pilot Zivi Nedivi and his copilot, the right wing of the Eagle was sheared off roughly two feet (60 cm) from the fuselage. The A-4 disintegrated and its pilot ejected and parachuted to safety, while the F-15 nosed down and entered a violent roll. Zivi decided to attempt recovery and engaged afterburner to increase speed, allowing him to regain control of the aircraft. The pilot was able to prevent stalling and maintain control because of the lift generated by the large horizontal surface area of the fuselage, the stabilators, and remaining wing areas. The F-15 landed at twice the normal speed to maintain the necessary descent and its tailhook was torn off during the landing. Zivi managed to bring his F-15 to a complete stop approximately 20 ft (6 m) from the end of the runway. He was later quoted as saying "It's highly likely that if I would have seen it clearly, I would have ejected..."; the fuel leak and vapors along the wing had prevented him from seeing what had happened to the wing itself. The aircraft was repaired and saw further combat service.
- On 19 March 1990, an F-15 from the 3rd Wing stationed at Elmendorf AFB, Alaska accidentally fired an AIM-9M Sidewinder missile at another F-15. The damaged aircraft was able to make an emergency landing; it was subsequently repaired and returned to service. This was not in combat, but does mark the first time an F-15 was struck by an air-to-air missile, accident or otherwise.
- On 22 November 1995, during air-intercept training over the Sea of Japan, a Japanese F-15J flown by Lt. Tatsumi Higuchi was shot down by an AIM-9L Sidewinder missile inadvertently fired by his wingman in an accident similar to the one that occurred on 19 March 1990. The pilot ejected safely. Both F-15Js involved were from JASDF 303rd Squadron, Komatsu AFB.
- On 26 March 2001, two US Air Force F-15Cs crashed near the summit of Ben Macdui in the Cairngorms during a low flying training exercise over the Scottish Highlands. Both Lieutenant Colonel Kenneth John Hyvonen and Captain Kirk Jones died in the accident, which resulted in a court martial for an RAF air traffic controller, who was later found not guilty.
- On 2 November 2007, a 27-year-old F-15C (AF Ser. No. 80-0034) of the 131st Fighter Wing, Missouri Air National Guard), crashed during air combat maneuvering training near St. Louis, Missouri. The pilot, Major Stephen W. Stilwell, ejected but suffered serious injuries. The crash was the result of an in-flight breakup due to structural failure. On 3 November 2007, all non-mission critical models of the F-15 were grounded pending the outcome of the crash investigation, and on the following day, grounded non-mission critical F-15s engaged in combat missions in the Middle East. By 13 November 2007, over 1,100 were grounded worldwide after Israel, Japan and Saudi Arabia grounded their aircraft as well. F-15Es were cleared on 15 November 2007 pending aircraft passing inspections. On 8 January 2008, the USAF cleared 60 percent of the F-15A/B/C/D fleet for return to flight. On 10 January 2008, the accident review board released its report stating the 2 November crash was related to the longeron not meeting drawing specifications. The Air Force cleared all its grounded F-15A-D fighters for flight on 15 February 2008 pending inspections, reviews and any needed repairs. In March 2008, Stilwell, the injured pilot, filed a lawsuit against Boeing, the F-15's manufacturer.
- On 20 February 2008, two F-15Cs from the 58th Fighter Squadron, 33rd Fighter Wing at Eglin AFB, Florida, flown by 1st Lt Ali Jivanjee and Capt Tucker Hamilton collided over the Gulf of Mexico during a training mission. Both pilots ejected and were rescued, but one died later from his injuries. The accident investigation report released 25 August 2008 found that the accident was the result of pilot error and not mechanical failure. Both pilots failed to clear their flight paths and anticipate their impending high-aspect, midair impact according to Brig Gen Joseph Reynes Jr., the leader of the investigation team.
Specifications (F-15C)
- Crew: 1: pilot
- Length: 63 ft 9 in (19.43 m)
- Wingspan: 42 ft 10 in (13.05 m)
- Height: 18 ft 6 in (5.63 m)
- Wing area: 608 ft² (56.5 m²)
- Airfoil: NACA 64A006.6 root, NACA 64A203 tip
- Empty weight: 28,000 lb (12,700 kg)
- Loaded weight: 44,500 lb (20,200 kg)
- Max. takeoff weight: 68,000 lb (30,845 kg)
- Powerplant: 2 × Pratt & Whitney F100-100 or −220 afterburning turbofans
- Fuel capacity: 13,455 lb (6,100 kg) internal
- Maximum speed:
- High altitude: Mach 2.5+ (1,650+ mph, 2,665+ km/h)
- Low altitude: Mach 1.2 (900 mph, 1,450 km/h)
- Combat radius: 1,061 nmi (1,222 mi, 1,967 km) for interdiction mission
- Ferry range: 3,450 mi (3,000 nmi, 5,550 km) with conformal fuel tanks and three external fuel tanks
- Service ceiling: 65,000 ft (20,000 m)
- Rate of climb: >50,000 ft/min (254 m/s)
- Wing loading: 73.1 lb/ft² (358 kg/m²)
- Thrust/weight: 1.12 (−220)
- Guns: 1× 20 mm (0.787 in) M61 Vulcan 6-barreled gatling cannon, 940 rounds
- Hardpoints: Total 11 (not including CFTs): two under-wing (each with additional two missile launch rails), four under-fuselage (for semi-recessed carriage of AIM-7 Sparrows) and a single centerline pylon station, optional fuselage pylons (which may include conformal fuel tanks, known initially as Fuel And Sensor Tactical (FAST) pack for use on the C model) with a capacity of 16,000 lb (7,300 kg) and provisions to carry combinations of:
- Northrop Grumman Electronic Systems AN/ALQ-131 electronic countermeasures pod
- Hazeltine AN/APX-76 or Raytheon AN/APX-119 Identify Friend/Foe (IFF) interrogator
- Magnavox AN/ALQ-128 Electronic Warfare Warning Set (EWWS) – part of Tactical Electronic Warfare Systems (TEWS)
- Loral AN/ALR-56 Radar warning receivers (RWR) – part of TEWS
- Northrop Grumman Electronic Systems ALQ-135 Internal Countermeasures System (ICS) – part of TEWS
- Marconi AN/ALE-45 Chaff/Flares dispenser system – part of TEWS
Aircraft on display
Although the F-15 continues to be a front-line fighter, a number of older USAF and IAF models have been retired, with several placed on outdoor display or in museums. These include:
United Kingdom
- 74-0131 - Wings of Liberty Memorial Park, RAF Lakenheath.
- 76-0020 - American Air Museum, Duxford.
United States
- 71-0280 - Lackland AFB, Texas.
- 71-0281 - Langley AFB, Virginia.
- 71-0283 - Defense Supply Center Richmond, Richmond, Virginia.
- 71-0285 - USAF Personnel Recruiting Office, St. Louis, Missouri.
- 71-0286- Octave Chanute Aerospace Museum, Rantoul, Illinois.
- 72-0119 "Streak Eagle" - in storage at the National Museum of the United States Air Force, Wright-Patterson AFB, Dayton, Ohio.
- 73-0085 - Museum of Aviation, Robins AFB, Warner Robins, Georgia.
- 73-0086 - Louisiana Military Museum, Jackson Barracks, New Orleans, Louisiana.
- 73-0099 (Marked as 77-0099) - Robins AFB, Warner Robins, Georgia.
- 74-0081 - Elmendorf AFB, Alaska.
- 74-0084 - Alaska Aviation Heritage Museum, Anchorage, Alaska.
- 74-0095 - Tyndall AFB, Panama City, Florida.
- 74-0114 - Mountain Home AFB, Idaho.
- 74-0117 - Langley AFB, Virginia.
- 74-0118 - Pima Air & Space Museum, Tucson, Arizona.
- 74-0119 - Castle Air Museum, Atwater, California.
- 74-0124 - Air Force Armament Museum, Eglin AFB, Florida.
- 75-0026 - National Warplane Museum, Elmira Corning Regional Airport, New York.
- 75-0045 - USS Alabama Battleship Memorial Park, Mobile, Alabama.
- 76-0008 - March Field Air Museum, Riverside, California.
- 76-0009 - Kingsley Field Air National Guard Base, Klamath Falls, Oregon.
- 76-0014 - Evergreen Aviation Museum, McMinnville, Oregon.
- 76-0024 - Peterson Air and Space Museum, Peterson AFB, Colorado.
- 76-0027 - National Museum of the United States Air Force, Wright-Patterson AFB, Dayton, Ohio.
- 76-0037 - Holloman AFB, New Mexico.
- 76-0040 - USAF Academy, Colorado Springs, Colorado.
- 76-0048 - McChord Air Museum, McChord AFB, Washington.
- 76-0063 - Pacific Aviation Museum, Ford Island, Joint Base Pearl Harbor-Hickam, Hawaii.
- 76-0066 - Portland Air National Guard Base, Oregon.
- 76-0076 (Marked as 85-0125) - roadside park, DeBary, Florida.
- 76-0080 - Jacksonville International Airport, Florida.
- 76-0088 - St. Louis International Airport, Lambert Field, Missouri.
- 76-0108 - Kelly AFB, Texas.
- 76-0110 - gate guard, Mountain Home AFB, Idaho.
- 77-0068 - Arnold AFB, Manchester, Tennessee.
- 77-0090 - Hill Aerospace Museum, Hill AFB, Utah.
- 77-0102 - Pacific Coast Air Museum, Charles M. Schulz-Sonoma County Airport, Santa Rosa, California. One of two Massachusetts Air National Guard 102d Fighter Wing aircraft scrambled in first response to terrorist air attacks on 11 September 2001.
- 77-0146 - Veterans Park, Callaway, Florida.
- 77-0150 - Yanks Air Museum, Chino, California.
- 73-0108 - Luke AFB, Arizona.
- 73-0114 - Air Force Flight Test Center Museum, Edwards AFB, California.
- 77-0161 - Seymour Johnson AFB, Goldsboro, North Carolina.
Notable appearances in media
The F-15 was the subject of the IMAX movie Fighter Pilot: Operation Red Flag, about the RED FLAG exercises. In Tom Clancy's nonfiction book, Fighter Wing (1995), a detailed analysis of the Air Force's premier fighter aircraft, the F-15 Eagle and its capabilities are showcased. Clancy's Red Storm Rising depicts the F-15 as an anti-satellite missile launcher.
The F-15 has also been a popular subject as a toy, and a fictional likeness of an aircraft similar to the F-15 has been used in cartoons, books, video games, animated television series, and animated films. The F-15 was mentioned in a veteran's old war story in the 2005 song Something to Be Proud Of by Montgomery Gentry.
See also
- Related development
- Aircraft of comparable role, configuration and era
- Related lists
- Although several F-15C aircraft were produced with APG-70 radar, all have been retrofitted to the AN/APG-63(V)1 configuration.
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- White, Josh. "2 F-15 Jets Crash: 1 Pilot Dies." Washington Post, 21 February 2008.
- Sirak, Michael with Marc Schanz. "Airman dies: Pilot error blamed." Air Force Magazine, Volume 91, Number 11, November 2008, p. 20.
- "F-15 Eagle fact sheet." U.S. Air Force, March 2008. Retrieved: 1 September 2011.
- Lambert 1993, p. 522.
- Davies 2002, Appendix 1.
- Spick 2000, p. 137.
- Parsch, Andreas. "AN/APG: Airborne fire control radars." Designation-Systems.Net, 20 November 2008. Retrieved: 27 September 2010.
- Schanz, Marc V., Assoc. Editor. "F-15s to Get New Radars." Aerospace World: Air Force Magazine, Journal of the Air Force Association Volume 90, Issue 6, p. 18, December 2007. ISSN 0730-6784.
- Parsch, Andreas. "AN/ALQ – Airborne Countermeasures Multipurpose/Special Equipment." Designation-systems.net, 9 October 2007. Retrieved: 27 September 2010.
- Parsch, Andreas. "AN/APX – Airborne Identification Radars." Designation-Systems.Net, 9 October 2007. Retrieved: 27 September 2010.
- Parsch, Andreas. "AN/ALR – Airborne Countermeasures Receivers." Designation-Systems.Net, 20 November 2008. Retrieved: 27 September 2010.
- Parsch, Andreas. "AN/ALE: Airborne countermeasures ejectors." Designation-Systems.Net, 20 November 2008. Retrieved: 27 September 2010.
- "F-15 Eagle/74-0085." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0109." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0083." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0088." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/73-0098." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/73-0107." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0131." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0020." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/71-0280." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/71-0281." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/71-0283." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/71-0285." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/71-0286." Octave Chanute Aerospace Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/72-0119." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/73-0085." Museum of Aviation.] Retrieved: 26 March 2013.
- "F-15 Eagle/73-0086." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/73-0099." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0081." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0084." Alaska Aviation Heritage Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0095." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0114." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0117." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0118." Pima Air & Space Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0119." Castle Air Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/74-0124." Air Force Armament Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/75-0026." National Warplane Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/75-0045." USS Alabama Battleship Memorial Park.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0008." March Field Air Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0009." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0014." Evergreen Aviation Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0024." Peterson Air and Space Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0027." National Museum of the USAF.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0037." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0040." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0048." McChord Air Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0063." Pacific Aviation Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0066." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0076." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0080." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0088." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0108." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/76-0110." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/77-0068." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/77-0090." Hill Aerospace Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/77-0102." Pacific Coast Air Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/77-0146." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/77-0150." Yanks Air Museum.] Retrieved: 26 March 2013.
- "F-15 Eagle/73-0108." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/73-0114." Warbird Registry.] Retrieved: 26 March 2013.
- "F-15 Eagle/77-0161." Warbird Registry.] Retrieved: 26 March 2013.
- Clancy, Tom. Fighter Wing: A Guided Tour of an Air Force Combat Wing. New York: Berkley Books, 1995. ISBN 0-425-14957-9.
- "Something To Be Proud Of: Lyrics." Cowboylyrics.com. Retrieved: 25 March 2012.
- Bowman, Martin W. US Military Aircraft. London: Bison Books, 1980. ISBN 0-89009-292-3.
- Davies, Steve. Boeing F-15E Strike Eagle, All-Weather Attack Aircraft. London: Airlife Publishing, Ltd., 2003. ISBN 1-84037-378-4.
- Davies, Steve. Combat Legend, F-15 Eagle and Strike Eagle. London: Airlife Publishing, Ltd., 2002. ISBN 1-84037-377-6.
- Davies, Steve and Doug Dildy. F-15 Eagle Engaged, The World's Most Successful Jet Fighter. Oxford, UK: Osprey Publishing, 2007. ISBN 978-1-84603-169-4.
- Eden, Paul and Soph Moeng, eds. The Complete Encyclopedia of World Aircraft. London: Amber Books Ltd., 2002. ISBN 0-7607-3432-1.
- Gething, Michael J. F-15 Eagle (Modern Fighting Aircraft). New York: Arco, 1983. ISBN 0-668-05902-8.
- Green, William and Gordon Swanborough. The Complete Book of Fighters. New York: Barnes & Noble Inc., 1988. ISBN 0-7607-0904-1.
- Gunston, Bill. American Warplanes. New York: Crescent Books. 1986. ISBN 0-517-61351-4.
- Huenecke, Klaus. Modern Combat Aircraft Design. Annapolis, Maryland: Naval Institute Press, 1987. ISBN 0-87021-426-8.
- Jenkins, Dennis R. F/A-18 Hornet: A Navy Success Story, pp. 1–8. New York: McGraw-Hill, 2000. ISBN 0-07-134696-1.
- Jenkins, Dennis R. McDonnell Douglas F-15 Eagle, Supreme Heavy-Weight Fighter. Hinckley, UK: Midland Publishing, 1998. ISBN 1-85780-081-8.
- Lambert, Mark, ed. Jane's All the World's Aircraft 1993–94. Alexandria, Virginia: Jane's Information Group Inc., 1993. ISBN 0-7106-1066-1.
- Scutts, Jerry. Supersonic Aircraft of USAF. New York: Mallard Press, 1989. ISBN 0-7924-5013-2.
- Spick, Mike, ed. The Great Book of Modern Warplanes. St. Paul, Minnesota: MBI, 2000. ISBN 0-7603-0893-4.
Further reading
- Braybrook, Roy. F-15 Eagle. London: Osprey Aerospace, 1991. ISBN 1-85532-149-1.
- Crickmore, Paul. McDonnell Douglas F-15 Eagle (Classic Warplanes series). New York: Smithmark Books, 1992. ISBN 0-8317-1408-5.
- Drendel, Lou. Eagle (Modern Military Aircraft Series). Carrollton, Texas: Squadron/Signal Publications, 1985. ISBN 0-8974-7271-3.
- Drendel, Lou and Don Carson. F-15 Eagle in action. Carrollton, Texas: Squadron/Signal Publications, 1976. ISBN 0-89747-023-0.
- Fitzsimons, Bernard. Modern Fighting Aircraft, F-15 Eagle. London: Salamander Books Ltd., 1983. ISBN 0-86101-182-1.
- Gething, Michael J. and Paul Crickmore. F-15 (Combat Aircraft series). New York: Crescent Books, 1992. ISBN 0-517-06734-X.
- Kinzey, Bert. The F-15 Eagle in Detail & Scale (Part 1, Series II). El Paso, Texas: Detail & Scale, Inc., 1978. ISBN 0-8168-5028-3.
- Rininger, Tyson V. F-15 Eagle at War. Minneapolis, USA: Zenith Press, 2009. ISBN 978-0-7603-3350-1.
|Wikimedia Commons has media related to: F-15 Eagle|
- F-15 Eagle USAF Fact Sheet
- F-15 Eagle history page on Boeing.com
- McDonnell Douglas F-15A, and F-15C on USAF National Museum web site
- F-15 Eagle in service with Israel
- F-15 page on GlobalSecurity.org
- The McDonnell Douglas F-15 Eagle page on Vectorsite.net |
Prophecy of the Popes
The Prophecy of the Popes (Latin: Prophetia Sancti Malachiae Archiepiscopi, de Summis Pontificibus) is a series of 112 short, cryptic phrases in Latin which purport to predict the Roman Catholic popes (along with a few antipopes), beginning with Pope Celestine II. The alleged prophecies were first published by Benedictine monk Arnold Wion in 1595. Wion attributes the prophecies to Saint Malachy, a 12th‑century Archbishop of Armagh, Ireland.
Given the very accurate description of popes up to 1590 and lack of accuracy after that year, historians generally conclude that the alleged prophecies are a fabrication written shortly before they were published. The Roman Catholic Church also dismisses them as forgery. The prophecies may have been created in an attempt to suggest that Cardinal Girolamo Simoncelli's bid for the papacy in the second conclave of 1590 was divinely ordained.
Many proponents of the prophecies claim that Pope Francis corresponds to "Peter the Roman" the pope described in the final prophecy, whose pontificate will allegedly bring the destruction of the city of Rome and usher in the beginning of the Apocalypse.
The alleged prophecies were first published in 1595 by a Benedictine named Arnold Wion in his Lignum Vitæ, a history of the Benedictine order. Wion attributed the prophecies to Saint Malachy, the 12th‑century Archbishop of Armagh. He explained that the prophecies had not, to his knowledge, ever been printed before, but that many were eager to see them. Wion includes both the alleged original prophecies, consisting of short, cryptic Latin phrases, as well as an interpretation applying the statements to historical popes up to Urban VII (pope for thirteen days in 1590), which Wion attributes to Alphonsus Ciacconius, an attribution which was refuted by Claude-François Menestrier in 1694.
According to an account put forward in 1871 by Abbé Cucherat, Malachy was summoned to Rome in 1139 by Pope Innocent II to receive two wool palliums for the metropolitan sees of Armagh and Cashel. While in Rome, Malachy purportedly experienced a vision of future popes, which he recorded as a sequence of cryptic phrases. This manuscript was then deposited in the Vatican Secret Archives, and forgotten about until its rediscovery in 1590, supposedly just in time for a papal conclave ongoing at the time.
Saint Bernard of Clairvaux, a contemporary biographer of Malachy who recorded the saint's alleged miracles, makes no mention of the prophecies, nor are they mentioned in any record prior to their 1595 publication.
Several historians have concluded that the prophecies are a late 16th‑century forgery. Spanish monk and scholar Benito Jerónimo Feijóo y Montenegro wrote in his Teatro Crítico Universal (1724–1739), in an entry called Purported prophecies, that the high level of accuracy of the alleged prophecies up until the date they were published, compared with their high level of inaccuracy after that date, is evidence that they were created around the time of publication. The prophecies and explanations given in Wion correspond very closely to a 1557 history of the popes by Onofrio Panvinio (including replication of errors made by Panvinio), which may indicate that the prophecies were written based on that source.
One theory to explain the creation of the prophecies, put forward by 17th-century French priest and encyclopaedist Louis Moréri, among others, is that they were spread by supporters of Cardinal Girolamo Simoncelli in support of his bid to become pope during the 1590 conclave to replace Urban VII. In the prophecies, the pope following Urban VII is given the description "Ex antiquitate Urbis" ("from the old city"), and Simoncelli was from Orvieto, which in Latin is Urbevetanum, old city. The prophecies may, therefore, have been created in an attempt to demonstrate that Simoncelli was destined to be pope. Simoncelli was not elected pope; Urban VII was succeeded by Pope Gregory XIV, born Niccolò Sfondrati.
The interpretation of the prophecies for pre-publication popes provided by Wion involves close correspondences between the mottos and the popes' birthplaces, family names, personal arms, and pre-papal titles. For example, the first motto, Ex castro Tiberis (from a castle on the Tiber), fits Pope Celestine II's birthplace in Città di Castello, on the Tiber.
Efforts to connect the prophecies to historical popes who were elected after its publication have been more strained. For example, Pope Clement XIII is referred to in a prophecy as Rosa Umbriae (the rose of Umbria), but was not from Umbria nor had any but the most marginal connection with the region, having been briefly pontifical governor of Rieti, at the time part of Umbria.
One writer notes that among the post-publication (post-1595) predictions there remain "some surprisingly appropriate phrases," while adding that "it is of course easy to exaggerate the list's accuracy by simply citing its successes," and that "other tags do not fit so neatly." Among the reported 'successes' are 'Religion depopulated' for Benedict XV (1914–22) whose papacy included World War One and the atheistic communist Russian Revolution; 'Light in the sky' for Leo XIII (1878-1903), with a comet in his coat of arms; and 'Flower of flowers' for Paul VI (1963–78), with fleur-de-lys in his coat of arms.
If we were to place the works of those who have repudiated the Prophecies of Malachy on scales and balance them against those who have accepted them, we would probably reach a fair equilibrium; however, the most important factor, namely the popularity of the prophecies, particularly among the ordinary people (as distinct from scholars), makes them as relevant to the second half of the twentieth century as they have ever been.
— Bander 1969, p. 10.
M.J. O'Brien, a Catholic priest who authored an 1880 monograph on the prophecies, provided a more critical assessment:
These prophecies have served no purpose. They are absolutely meaningless. The Latin is bad. It is impossible to attribute such absurd triflings... to any holy source. Those who have written in defence of the prophecy... have brought forward scarcely an argument in their favour. Their attempts at explaining the prophecies after 1590 are, I say with all respect, the sorriest trifling.
— O'Brien 1880, p. 110.
In recent times, some interpreters of prophetic literature have drawn attention to the prophecies due to their imminent conclusion; if the list of descriptions is matched on a one-to-one basis to the list of historic popes since the prophecies' publication, Benedict XVI (2005-2013) would correspond to the second to last of the papal descriptions, Gloria olivae (the glory of the olive).
In persecutione extrema S.R.E. sedebit.
Petrus Romanus, qui pascet oves in multis tribulationibus, quibus transactis civitas septicollis diruetur, & judex tremendus judicabit populum suum. Finis.
This may be translated into English as:
In the final persecution of the Holy Roman Church, there will sit [i.e., as bishop].
Peter the Roman, who will pasture his sheep in many tribulations, and when these things are finished, the city of seven hills [i.e. Rome] will be destroyed, and the dreadful judge will judge his people. The End.
Several historians and interpreters of the prophecies note that they leave open the possibility of unlisted popes between "the glory of the olive" and the final pope, "Peter the Roman". In the Lignum Vitae, the line In persecutione extrema S.R.E. sedebit. forms a separate sentence and paragraph of its own. While often read as part of the "Peter the Roman" prophecy, other interpreters view it as a separate, incomplete sentence explicitly referring to additional popes between "the glory of the olive" and "Peter the Roman".
Popes and corresponding mottos
The list can be divided into two groups; one of the 74 popes and antipopes who reigned prior to the appearance of the prophecies c. 1590, for whom the connection between the motto and the pope is consistently clear. The other is of the 38 mottos attributed to popes who have reigned since 1590, for whom the connection between the motto and the pope is often strained or totally absent and could be viewed as shoehorning or postdiction.
René Thibaut divides the table at a different point, between the 71st and 72nd motto, asserting that there is a change in style at this point. He uses this distinction to put forward the view that the first 71 mottos are post-dated forgeries, while the remainder are genuine. Hildebrand Troll echoes this view, noting that mottos 72-112 use a symbolic language related to the character of the pope and his papacy, in contrast to the more literal mottos for earlier popes.
Popes and antipopes 1143–1590 (pre-publication)
The text on the silver lines below reproduces the original text (including punctuation and orthography) of the 1595 Lignum Vitae, which consisted of three parallel columns for the popes before 1590. The first column contained the motto, the second the name of the pope or antipope to whom it was attached (with occasional errors), and the third an explanation of the motto. There are some indications that both the mottos and explanations were the work of a single 16th century individual. The original list was unnumbered.
|Pre-appearance Popes (1143–1590)|
|Motto No.||Motto (Translation)||Regnal Name (Reign)||Name||Explanation Provided in Lignum Vitae||Coat of Arms|
|Ex caſtro Tiberis.||Cœleſtinus. ij.||Typhernas.|
|1.||From a castle of the Tiber||Celestine II (1143–1144)||Guido de Castello||An inhabitant of Tifernum.
Celestine II was born in Città di Castello (formerly called Tifernum-Tiberinum), on the banks of the Tiber.
|Inimicus expulſus.||Lucius. ij.||De familia Caccianemica.|
|2.||Enemy expelled||Lucius II (1144–1145)||Gherardo Caccianemici del Orso||Of the Caccianemici family.
According to Wion, this motto refers to Lucius II's family name, Caccianemici; in Italian, “Cacciare” means “to drive out” and “nemici” means “enemies”. While he has been traditionally viewed as being part of this family, it is doubtful whether he actually was; moreover, even if he actually belonged to that family, the attribution of the surname Caccianemici is certainly anachronistic.
|Ex magnitudine mõtis.||Eugenius. iij.||Patria Ethruſcus oppido Montis magni.|
|3.||From the great mountain||Eugene III (1145–1153)||Bernardo dei Paganelli di Montemagno||Tuscan by nation, from the town of Montemagno.
According to Wion, the motto refers to Eugene III’s birthplace, “Montemagno”, a village near Pisa. But according to other sources he was born in Pisa in modest family.
|Abbas Suburranus.||Anaſtaſius. iiij.||De familia Suburra.|
|4.||Abbot from Subbura||Anastasius IV (1153–1154)||Corrado di Suburra||From the Suburra family. He was traditionally referred to as abbot of the canon regulars of St. Ruf in Avignon, but modern scholars have established that he actually belonged to the secular clergy.|
|De rure albo.||Adrianus. iiij.||Vilis natus in oppido Sancti Albani.|
|5.||From the white countryside||Adrian IV (1154–1159)||Nicholas Breakspear||Humbly born in the town of St. Albans.
Most likely a reference to Adrian IV's birthplace near St Albans, Hertfordshire.
|Ex tetro carcere.||Victor. iiij.||Fuit Cardinalis S. Nicolai in carcere Tulliano.|
|6.||Out of a loathsome prison.||Victor IV, Antipope (1159–1164)||Ottaviano Monticello||He was a cardinal of St. Nicholas in the Tullian prison.
Victor IV may have held the title San Nicola in Carcere.
|Via Tranſtiberina.||Calliſtus. iij. [sic]||Guido Cremenſis Cardinalis S. Mariæ Tranſtiberim.|
|7.||Road across the Tiber.||Callixtus III, Antipope (1168–1178)||Giovanni di Strumi||Guido of Crema, Cardinal of St. Mary across the Tiber.
Wion reverses the names and order of Antipopes Callixtus III (John of Struma) and Paschal III (Guido of Crema). Paschal, not Callixtus, was born Guido of Crema and held the title of Santa Maria in Trastevere, to which the motto applies.
|De Pannonia Thuſciæ.||Paſchalis. iij. [sic]||Antipapa. Hungarus natione, Epiſcopus Card. Tuſculanus.|
|8.||From Tusculan Hungary.||Paschal III, Antipope (1164–1168)||Guido di Crema||Antipope. A Hungarian by birth, Cardinal Bishop of Tusculum.
As noted above, this motto applies not to Paschal III, but to Callixtus III, who allegedly was Hungarian. However, Callixtus was Cardinal Bishop of Albano, not of Tusculum.
|Ex anſere cuſtode.||Alexander. iij.||De familia Paparona.|
|9.||From the guardian goose||Alexander III (1159–1181)||Rolando (or Orlando) of Siena||Of the Paparoni family.
Alexander III may have been from the Bandinella family, which was afterwards known as the Paparona family, which featured a goose on its coat of arms. There is debate whether Alexander III was in fact of that family.
|Lux in oſtio.||Lucius. iij.||Lucenſis Card. Oſtienſis.|
|10.||A light in the door||Lucius III (1181–1185)||Ubaldo Allucingoli||A Luccan Cardinal of Ostia.
The motto is a wordplay on "Lucius" or "Lucca" and "Ostia".
|Sus in cribro.||Vrbanus. iij.||Mediolanenſis, familia cribella, quæ Suem pro armis gerit.|
|11.||Pig in a sieve||Urban III (1185–1187)||Umberto Crivelli||A Milanese, of the Cribella (Crivelli) family, which bears a pig for arms.
Urban III's family name Crivelli means "a sieve" in Italian; his arms included a sieve and two pigs.
|Enſis Laurentii.||Gregorius. viij.||Card. S. Laurentii in Lucina, cuius inſignia enſes falcati.|
|12.||The sword of Lawrence||Gregory VIII (1187)||Alberto De Morra||Cardinal of St. Lawrence in Lucina, of whom the arms were curved swords.
Gregory VIII was Cardinal of St. Lawrence and his arms featured crossed swords.
|De Schola exiet.||Clemens. iij.||Romanus, domo Scholari.|
|13.||He will come from school||Clement III (1187–1191)||Paolo Scolari||A Roman, of the house of Scolari.
The motto is a play on words on Clement III's surname.
|De rure bouenſi.||Cœleſtinus. iij.||Familia Bouenſi.|
|14.||From cattle country||Celestine III (1191–1198)||Giacinto Bobone||Bovensis family.
The reference to cattle is a wordplay on Celestine III's surname, Bobone.
|Comes Signatus.||Innocentius. iij.||Familia Comitum Signiæ.|
|15.||Designated count||Innocent III (1198–1216)||Lotario dei Conti di Segni||Family of the Counts of Signia (Segni)
The motto is a direct reference to Innocent III's family name.
|Canonicus de latere.||Honorius. iij.||Familia Sabella, Canonicus S. Ioannis Lateranensis.|
|16.||Canon from the side||Honorius III (1216–1227)||Cencio Savelli||Savelli family, canon of St. John Lateran
The claim in Wion that Honorius III was a canon of St. John Lateran is contested by some historians.
|Auis Oſtienſis.||Gregorius. ix.||Familia Comitum Signiæ Epiſcopus Card. Oſtienſis.|
|17.||Bird of Ostia||Gregory IX (1227–1241)||Ugolino dei Conti di Segni||Family of the Counts of Segni, Cardinal Bishop of Ostia.
Before his election to the papacy, Ugolino dei Conti was the Cardinal Bishop of Ostia, and his coat of arms depict an eagle.
|Leo Sabinus.||Cœleſtinus iiij.||Mediolanenſis, cuius inſignia Leo, Epiſcopus Card. Sabinus.|
|18.||Sabine Lion||Celestine IV (1241)||Goffredo Castiglioni||A Milanese, whose arms were a lion, Cardinal Bishop of Sabina.
Celestine IV was Cardinal Bishop of Sabina and his armorial bearing had a lion in it.
|Comes Laurentius.||Innocentius iiij.||domo flisca, Comes Lauaniæ, Cardinalis S. Laurentii in Lucina.|
|19.||Count Lawrence||Innocent IV (1243–1254)||Sinibaldo Fieschi||Of the house of Flisca (Fieschi), Count of Lavagna, Cardinal of St. Lawrence in Lucina.
The motto, as explained in Wion, is a reference to Innocent IV's father, the Count of Lavagna, and his title Cardinal of St. Lawrence in Lucina.
|Signum Oſtienſe.||Alexander iiij.||De comitibus Signiæ, Epiſcopus Card. Oſtienſis.|
|20.||Sign of Ostia||Alexander IV (1254–1261)||Renaldo dei Signori di Ienne||Of the counts of Segni, Cardinal Bishop of Ostia.
The motto refers to Alexander IV's being Cardinal Bishop of Ostia and member of the Conti-Segni family.
|Hieruſalem Campanię.||Vrbanus iiii.||Gallus, Trecenſis in Campania, Patriarcha Hieruſalem.|
|21.||Jerusalem of Champagne||Urban IV (1261–1264)||Jacques Pantaleon||A Frenchman, of Trecae (Troyes) in Champagne, Patriarch of Jerusalem.
The motto refers to Urban IV's birthplace of Troyes, Champagne, and title Patriarch of Jerusalem.
|Draco depreſſus.||Clemens iiii.||cuius inſignia Aquila vnguibus Draconem tenens.|
|22.||Dragon pressed down||Clement IV (1265–1268)||Guido Fulcodi||Whose badge is an eagle holding a dragon in his talons.
According some sources, Clement IV's coat of arms depicted an eagle clawing a dragon. Other sources indicate that it was instead six fleurs-de-lis.
|Anguinus uir.||Gregorius. x.||Mediolanenſis, Familia vicecomitum, quæ anguẽ pro inſigni gerit.|
|23.||Snaky man||Gregory X (1271–1276)||Teobaldo Visconti||A Milanese, of the family of Viscounts (Visconti), which bears a snake for arms.
The Visconti coat of arms had a large serpent devouring a male child feet first; sources conflict as to whether Gregory X used this for his papal arms.
|Concionator Gallus.||Innocentius. v.||Gallus, ordinis Prædicatorum.|
|24.||French Preacher||Innocent V (1276)||Pierre de Tarentaise||A Frenchman, of the Order of Preachers.
Innocent V was born in what is now south-eastern France and was a member of the order of Preachers.
|Bonus Comes.||Adrianus. v.||Ottobonus familia Fliſca ex comitibus Lauaniæ.|
|25.||Good Count||Adrian V (1276)||Ottobono Fieschi||Ottobono, of the Fieschi family, from the counts of Lavagna.
The Fieschi family were counts of Lavagna and a wordplay on "good" can be made with Adrian V's first name, Ottobono.
|Piſcator Thuſcus.||Ioannes. xxi.||antea Ioannes Petrus Epiſcopus Card. Tuſculanus.|
|26.||Tuscan Fisherman||John XXI (1276–1277)||Pedro Julião||Formerly John Peter, Cardinal Bishop of Tusculum.
John XXI had been the Cardinal Bishop of Tusculum, and shared his first name with Saint Peter, a fisherman.
|Roſa compoſita.||Nicolaus. iii.||Familia Vrſina, quæ roſam in inſigni gerit, dictus compoſitus.|
|27.||Composite Rose||Nicholas III (1277–1280)||Giovanni Gaetano Orsini||Of the Ursina (Orsini) family, which bears a rose on its arms, called 'composite'.
Nicholas III bore a rose in his coat of arms.
|Ex teloneo liliacei Martini.||Martinus. iiii.||cuius inſignia lilia, canonicus, & theſaurarius S. Martini Turonen[sis].|
|28.||From the tollhouse of Martin of the lilies||Martin IV (1281–1285)||Simone de Brion||Whose arms were lilies, canon and treasurer of St. Martin of Tours.
Martin IV was Canon and Treasurer at the Church of St. Martin in Tours, France. Wion's assertion that his arms featured lilies is incorrect.
|Ex roſa leonina.||Honorius. iiii.||Familia Sabella inſignia roſa à leonibus geſtata.|
|29.||Out of the leonine rose||Honorius IV (1285–1287)||Giacomo Savelli||Of the Sabella (Savelli) family, arms were a rose carried by lions.
Honorius IV's coat of arms was emblazoned with two lions supporting a rose.
|Picus inter eſcas.||Nicolaus. iiii.||Picenus patria Eſculanus.|
|30.||Woodpecker between food||Nicholas IV (1288–1292)||Girolamo Masci||A Picene by nation, of Asculum (Ascoli).
The motto is likely an obscure wordplay on Nicholas IV's birthplace in Ascoli, in Picenum.
|Ex eremo celſus.||Cœleſtinus. v.||Vocatus Petrus de morrone Eremita.|
|31.||Raised out of the desert||St. Celestine V (1294)||Pietro Di Murrone||Called Peter de Morrone, a hermit.
Prior to his election, Celestine V was a hermit (eremita, literally a dweller in the eremus, or desert).
|Ex undarũ bn̑dictione.||Bonifacius. viii.||Vocatus prius Benedictus, Caetanus, cuius inſignia undæ.|
|32.||From the blessing of the waves||Boniface VIII (1294–1303)||Benedetto Caetani||Previously called Benedict, of Gaeta, whose arms were waves.
Boniface VIII's coat of arms had a wave through it. Also a play on words, referring to the pope's Christian name, "Benedetto."
|Concionator patereus. [sic]||Benedictus. xi.||qui uocabatur Frater Nicolaus, ordinis Prædicatorum.|
|33.||Preacher From Patara||Benedict XI (1303–1304)||Nicholas Boccasini||Who was called Brother Nicholas, of the order of Preachers.
Benedict XI belonged to the Order of Preachers, and his namesake Saint Nicholas was from Patara. O'Brien notes, "Everything leads us to suspect that the author and interpreter of the prophecy is one and the same person. The pretended interpreter who knew that Patare was the birthplace of St. Nicholas forgot that others may not be aware of the fact, and that therefore the explanation would be thrown away on them."
|De feſſis aquitanicis.||Clemens V.||natione aquitanus, cuius inſignia feſſæ erant.|
|34.||From the fesses of Aquitaine||Clement V (1305–1314)||Bertrand de Got||An Aquitanian by birth, whose arms were fesses.
Clement V was Bishop of St-Bertrand-de-Comminges in Aquitaine, and eventually became Archbishop of Bordeaux, also in Aquitaine. His coat of arms displays three horizontal bars, known in heraldry as fesses.
|De ſutore oſſeo.||Ioannes XXII.||Gallus, familia Oſſa, Sutoris filius.|
|35.||From a bony cobbler||John XXII (1316–1334)||Jacques Duese||A Frenchman, of the Ossa family, son of a cobbler.
John XXII's family name was Duèze or D'Euse, the last of which might be back-translated into Latin as Ossa ("bones"), the name Wion gives. The popular legend that his father was a cobbler is dubious.
|Coruus ſchiſmaticus.||Nicolaus V.||qui uocabatur F. Petrus de corbario, contra Ioannem XXII. Antipapa Minorita.|
|36.||Schismatic crow||Nicholas V, Antipope (1328–1330)||Pietro Rainalducci di Corvaro||Who was called Brother Peter of Corbarium (Corvaro), the Minorite antipope opposing John XXII.
The motto is a play on words, referring to Pietro di Corvaro's last name.
|Frigidus Abbas.||Benedictus XII.||Abbas Monaſterii fontis frigidi.|
|37.||Cold abbot||Benedict XII (1334–1342)||Jacques Fournier||Abbot of the monastery of the cold spring.
Benedict XII was an abbot in the monastery of Fontfroide ("cold spring").
|De roſa Attrebatenſi.||Clemens VI.||Epiſcopus Attrebatenſis, cuius inſignia Roſæ.|
|38.||From the rose of Arras||Clement VI (1342–1352)||Pierre Roger||Bishop of Arras, whose arms were roses.
Clement VI was Bishop of Arras (in Latin, Episcopus Attrebatensis) and his armorial bearings were emblazoned with six roses.
|De mõtibus Pãmachii.||Innocentius VI.||Cardinalis SS. Ioannis & Pauli. T. Panmachii, cuius inſignia ſex montes erant.|
|39.||From the mountains of Pammachius||Innocent VI (1352–1362)||Etienne Aubert||Cardinal of Saints John and Paul, Titulus of Pammachius, whose arms were six mountains.
Innocent VI was Cardinal Priest of Pammachius. Wion and Panvinio describe his arms as depicting six mountains, though other sources do not.
|Gallus Vicecomes.||Vrbanus V.||nuncius Apoſtolicus ad Vicecomites Mediolanenſes.|
|40.||French viscount||Urban V (1362–1370)||Guglielmo De Grimoard||Apostolic nuncio to the Viscounts of Milan.
Urban V was French. Wion indicates he was Apostolic Nuncio to the Viscounts of Milan.
|Nouus de uirgine forti.||Gregorius XI.||qui uocabatur Petrus Belfortis, Cardinalis S. Mariæ nouæ.|
|41.||New man from the strong virgin||Gregory XI (1370–1378)||Pierre Roger de Beaufort||Who was called Peter Belfortis (Beaufort), Cardinal of New St. Mary's.
The motto refers to Gregory XI's surname and his title Cardinal of Santa Maria Nuova.
|Decruce Apoſtolica. [sic]||Clemens VII.||qui fuit Preſbyter Cardinalis SS. XII. Apoſtolorũ cuius inſignia Crux.|
|42.||From the apostolic cross||Clement VII, Antipope (1378–1394)||Robert, Count of Geneva||Who was Cardinal Priest of the Twelve Holy Apostles, whose arms were a cross.
Clement VII's coat of arms showed a cross and he held the title Cardinal Priest of the Twelve Holy Apostles.
|Luna Coſmedina.||Benedictus XIII.||antea Petrus de Luna, Diaconus Cardinalis S. Mariæ in Coſmedin.|
|43.||Cosmedine moon.||Benedict XIII, Antipope (1394–1423)||Peter de Luna||Formerly Peter de Luna, Cardinal Deacon of St. Mary in Cosmedin.
The motto refers to Benedict XIII's surname and title.
|Schiſma Barchinoniũ.||Clemens VIII.||Antipapa, qui fuit Canonicus Barchinonenſis.|
|44.||Schism of the Barcelonas||Clement VIII, Antipope (1423–1429)||Gil Sanchez Muñoz||Antipope, who was a canon of Barcelona.|
|De inferno prægnãti.||Vrbanus VI.||Neapolitanus Pregnanus, natus in loco quæ dicitur Infernus.|
|45.||From a pregnant hell.||Urban VI (1378–1389)||Bartolomeo Prignano||The Neapolitan Prignano, born in a place which is called Inferno.
Urban VI's family name was Prignano or Prignani, and he was native to a place called Inferno near Naples.
|Cubus de mixtione.||Bonifacius. IX.||familia tomacella à Genua Liguriæ orta, cuius inſignia Cubi.|
|46.||Square of mixture||Boniface IX (1389–1404)||Pietro Tomacelli||Of the Tomacelli family, born in Genoa in Liguria, whose arms were cubes.
Boniface IX's coat of arms includes a bend checky — a wide stripe with a checkerboard pattern.
|De meliore ſydere.||Innocentius. VII.||uocatus Coſmatus de melioratis Sulmonenſis, cuius inſignia ſydus.|
|47.||From a better star||Innocent VII (1404–1406)||Cosmo Migliorati||Called Cosmato dei Migliorati of Sulmo, whose arms were a star.
The motto is a play on words, "better" (melior) referring to Innocent VII's last name, Migliorati (Meliorati). There is a shooting star on his coat of arms.
|Nauta de Ponte nigro.||Gregorius XII.||Venetus, commendatarius eccleſiæ Nigropontis.|
|48.||Sailor from a black bridge||Gregory XII (1406–1415)||Angelo Correr||A Venetian, commendatary of the church of Negroponte.
Gregory XII was born in Venice (hence mariner) and was commendatary of Chalkis, then called Negropont.
|Flagellum ſolis.||Alexander. V.||Græcus Archiepiſcopus Mediolanenſis, inſignia Sol.|
|49.||Whip of the sun||Alexander V, Antipope (1409–1410)||Petros Philarges||A Greek, Archbishop of Milan, whose arms were a sun.
Alexander V's coat of arms featured a sun, the wavy rays may explain the reference to a whip.
|Ceruus Sirenæ.||Ioannes XXIII.||Diaconus Cardinalis S. Euſtachii, qui cum ceruo depingitur, Bononiæ legatus, Neapolitanus.|
|50.||Stag of the siren||John XXIII, Antipope (1410–1415)||Baldassarre Cossa||Cardinal Deacon of St. Eustace, who is depicted with a stag; legate of Bologna, a Neapolitan.
John XXIII was a cardinal with the title of St. Eustachius, whose emblem is a stag, and was originally from Naples, which has the emblem of the siren.
|Corona ueli aurei.||Martinus V.||familia colonna, Diaconus Cardinalis S. Georgii ad uelum aureum.|
|51.||Crown of the golden curtain||Martin V (1417–1431)||Oddone Colonna||Of the Colonna family, Cardinal Deacon of St. George at the golden curtain.
The motto is a reference to Martin V's family name and cardinal title of San Giorgio in Velabro.
|Lupa Cœleſtina,||Eugenius. IIII.||Venetus, canonicus antea regularis Cœleſtinus, & Epiſcopus Senẽſis.|
|52.||Heavenly she-wolf||Eugene IV (1431–1447)||Gabriele Condulmaro||A Venetian, formerly a regular Celestine canon, and Bishop of Siena.
Eugene IV belonged to the order of the Celestines and was the Bishop of Siena which bears a she-wolf on its arms.
|Amator Crucis.||Felix. V.||qui uocabatur Amadæus Dux Sabaudiæ, inſignia Crux.|
|53.||Lover of the cross||Felix V, Antipope (1439–1449)||Amadeus, Duke of Savoy||Who was called Amadeus, Duke of Savoy, arms were a cross.
The motto is a reference to Felix V's given name, Amadeus, and arms, which featured the cross of Savoy.
|De modicitate Lunæ.||Nicolaus V.||Lunenſis de Sarzana, humilibus parentibus natus.|
|54.||From the meanness of Luna||Nicholas V (1447–1455)||Tommaso Parentucelli||A Lunese of Sarzana, born to humble parents.
Nicholas V was born in the diocese of Luni, the ancient name of which was Luna.
|Bos paſcens.||Calliſtus. III.||Hiſpanus, cuius inſignia Bos paſcens.|
|55.||Pasturing ox||Callixtus III (1455–1458)||Alfonso Borja||A Spaniard, whose arms were a pasturing ox.
Callixtus III's coat of arms featured an ox.
|De Capra & Albergo.||Pius. II.||Senenſis, qui fuit à Secretis Cardinalibus Capranico & Albergato.|
|56.||From a nanny-goat and an inn||Pius II (1458–1464)||Enea Silvio de Piccolomini||A Sienese, who was secretary to Cardinals Capranicus and Albergatus.
Pius II was secretary to Cardinal Domenico Capranica and Cardinal Albergatti before he was elected Pope.
|De Ceruo & Leone.||Paulus. II.||Venetus, qui fuit Commendatarius eccleſiæ Ceruienſis, & Cardinalis tituli S. Marci.|
|57.||From a stag and lion||Paul II (1464–1471)||Pietro Barbo||A Venetian, who was commendatary of the church of Cervia, and Cardinal of the title of St. Mark.
The motto refers to his Bishopric of Cervia (punning on cervus, "a stag") and his Cardinal title of St. Mark (symbolized by a winged lion).
|Piſcator minorita.||Sixtus. IIII.||Piſcatoris filius, Franciſcanus.|
|58.||Minorite fisherman||Sixtus IV (1471–1484)||Francesco Della Rovere||Son of a fisherman, Franciscan.
Sixtus IV was born the son of a fisherman and a member of the Franciscans, also known as "Minorites" (which was founded in 1209, after Malachy's death.)
|Præcurſor Siciliæ.||Innocentius VIII.||qui uocabatur Ioãnes Baptiſta, & uixit in curia Alfonſi regis Siciliæ.|
|59.||Precursor of Sicily||Innocent VIII (1484–1492)||Giovanni Battista Cibò||Who was called John Baptist, and lived in the court of Alfonso, king of Sicily.
Innocent VIII was from Sicily. "Precursor" may be explained as an allusion to his birth name, after John the Baptist, the precursor of Christ.
|Bos Albanus in portu.||Alexander VI.||Epiſcopus Cardinalis Albanus & Portuenſis, cuius inſignia Bos.|
|60.||Bull of Alba in the harbor||Alexander VI (1492–1503)||Rodrigo de Borgia||Cardinal Bishop of Albano and Porto, whose arms were a bull.
In 1456, he was made a Cardinal and he held the titles of Cardinal Bishop of Albano and Porto, and his arms featured an ox.
|De paruo homine.||Pius. III.||Senenſis, familia piccolominea.|
|61.||From a small man||Pius III (1503)||Francesco Todeschini Piccolomini||A Sienese, of the Piccolomini family.
Pius III's family name was Piccolomini, from piccolo "small" and uomo "man".
|Fructus Iouis iuuabit.||Iulius. II.||Ligur, eius inſignia Quercus, Iouis arbor.|
|62.||The fruit of Jupiter will help||Julius II (1503–1513)||Giuliano Della Rovere||A Genoese, his arms were an oak, Jupiter's tree.
On Julius II's arms was an oak tree, which was sacred to Jupiter.
|De craticula Politiana.||Leo. X.||filius Laurentii medicei, & ſcholaris Angeli Politiani.|
|63.||From a Politian gridiron||Leo X (1513–1521)||Giovanni de Medici||Son of Lorenzo de' Medici, and student of Angelo Poliziano.
Leo X's educator and mentor was Angelo Poliziano. The “Gridiron” in the motto evidently refers to St. Lawrence, who was martyred on a gridiron. This is a rather elliptical allusion to Lorenzo the Magnificent, who was Giovanni’s father.
|Leo Florentius.||Adrian. VI.||Florẽtii filius, eius inſignia Leo.|
|64.||Florentian lion||Adrian VI (1522–1523)||Adriaen Florenszoon Boeyens||Son of Florentius, his arms were a lion.
Adrian VI's coat of arms had two lions on it, and his name is sometimes given as Adrian Florens, or other variants, from his father's first name Florens (Florentius).
|Flos pilei ægri.||Clemens. VII.||Florentinus de domo medicea, eius inſignia pila, & lilia.|
|65.||Flower of the sick man's pill||Clement VII (1523–1534)||Giulio de Medici||A Florentine of the Medicean house, his arms were pill-balls and lilies.
The Medici coat of arms was emblazoned with six medical balls. One of these balls, the largest of the six, was emblazoned with the Florentine lily.
|Hiacinthus medicorũ.||Paulus. III.||Farneſius, qui lilia pro inſignibus geſtat, & Card. fuit SS. Coſme, & Damiani.|
|66.||Hyacinth of the physicians||Paul III (1534–1549)||Alessandro Farnese||Farnese, who bore lilies for arms, and was Cardinal of Saints Cosmas and Damian.
According to some sources, Paul III's coat of arms were charged with hyacinths, and he was cardinal of Saints Cosmas and Damian, both doctors.
|De corona montana.||Iulius. III.||antea uocatus Ioannes Maria de monte.|
|67.||From the mountainous crown||Julius III (1550–1555)||Giovanni Maria Ciocchi del Monte||Formerly called Giovanni Maria of the Mountain (de Monte)
His coat of arms showed mountains and laurel crowns (chaplets).
|Frumentum flocidum. [sic]||Marcellus. II.||cuius inſignia ceruus & frumẽtum, ideo floccidum, quod pauco tempore uixit in papatu.|
|68.||Trifling grain||Marcellus II (1555)||Marcello Cervini||Whose arms were a stag and grain; 'trifling', because he lived only a short time as pope.
His coat of arms showed a stag and ears of wheat.
|De fide Petri.||Paulus. IIII.||antea uocatus Ioannes Petrus Caraffa.|
|69.||From Peter's faith||Paul IV (1555–1559)||Giovanni Pietro Caraffa||Formerly called John Peter Caraffa.
Paul IV is said to have used his second Christian name Pietro.
|Eſculapii pharmacum.||Pius. IIII.||antea dictus Io. Angelus Medices.|
|70.||Aesculapius' medicine||Pius IV (1559–1565)||Giovanni Angelo de Medici||Formerly called Giovanni Angelo Medici.
The motto is likely a simple allusion to Pius IV's family name.
|Angelus nemoroſus.||Pius. V.||Michael uocatus, natus in oppido Boſchi.|
|71.||Angel of the grove||St. Pius V (1566–1572)||Antonio Michele Ghisleri||Called Michael, born in the town of Bosco.
Pius V was born in Bosco, Lombardy; the placename means grove. His name was 'Antonio Michele Ghisleri', and Michele relates to the archangel. O'Brien notes here that many of the prophecies contain plays on Italian words, which are not made explicit in the explanations provided in the Lignum Vitae.
|Medium corpus pilarũ.||Gregorius. XIII.||cuius inſignia medius Draco, Cardinalis creatus à Pio. IIII. qui pila in armis geſtabat.|
|72.||Half body of the balls||Gregory XIII (1572–1585)||Ugo Boncompagni||Whose arms were a half-dragon; a Cardinal created by Pius IV who bore balls in his arms.
The "balls" in the motto refer to Pope Pius IV, who had made Gregory a cardinal. Pope Gregory had a dragon on his coat of arms with half a body.
|Axis in medietate ſigni.||Sixtus. V.||qui axem in medio Leonis in armis geſtat.|
|73.||Axle in the midst of a sign.||Sixtus V (1585–1590)||Felice Peretti||Who bears in his arms an axle in the middle of a lion.
This is a rather straightforward description of the Sixtus V's coat of arms.
|De rore cœli.||Vrbanus. VII.||qui fuit Archiepiſcopus Roſſanenſis in Calabria, ubi mãna colligitur.|
|74.||From the dew of the sky||Urban VII (1590)||Giovanni Battista Castagna||Who was Archbishop of Rossano in Calabria, where manna is collected.
He had been Archbishop of Rossano in Calabria where sap called "the dew of heaven" is gathered from trees.
Popes 1590 to present (post-publication)
For this group of popes, the published text only provides names for the first three (i.e., those who were popes between the appearance of the text c. 1590, and its publication in 1595) and provides no explanations.
|Post-appearance Popes (1590–present)|
|Motto No.||Motto (Translation)||Regnal Name (Reign)||Name||Interpretations and Criticisms||Coat of Arms|
|Ex antiquitate Vrbis.||Gregorius. XIIII.|
|75.||Of the antiquity of the city / From the old city||Gregory XIV (1590–1591)||Niccolo Sfondrati||This may have been intended by the author of the prophecies to suggest that Cardinal Girolamo Simoncelli was destined to succeed Urban VII. Simoncelli was from Orvieto, which in Latin is Urbs vetus, old city. Simoncelli was not elected pope, however, Niccolo Sfondrati was, who took the name Gregory XIV. Proponents of the prophecies have attempted to explain it by noting that Gregory XIV's father was a senator of the ancient city of Milan, and the word "senator" is derived from the Latin senex, meaning old man, or that Milan is the "old city" in question, having been founded c. 400 BCE.|
|Pia ciuitas in bello.||Innocentius. IX.|
|76.||Pious citizens in war||Innocent IX (1591)||Giovanni Antonio Facchinetti||Proponents of the prophecies have suggested different interpretations to relate this motto to Innocent IX, including references to his birthplace of Bologna or title of Patriarch of Jerusalem.|
|Crux Romulea.||Clemens. VIII.|
|77.||Cross of Romulus||Clement VIII (1592–1605)||Ippolito Aldobrandini||Proponents of the prophecies have suggested different interpretations to relate this motto to Clement VIII, including linking it to the embattled bend on his arms or the war between Catholic Ireland and Protestant England during his papacy.|
|78.||Wavy man||Leo XI (1605)||Alessandro Ottaviano De Medici||This may have been intended by the author of the prophecies to suggest to his audience a possible heraldic design, but it does not correspond to Leo XI's Medici arms. Proponents of the prophecies have suggested different interpretations to relate this motto to this pope, including relating it to his short reign "passing like a wave."|
|79.||Wicked race||Paul V (1605–1621)||Camillo Borghese||Proponents of the prophecies have suggested it is a reference to the dragon and the eagle on Paul V's arms.|
|In tribulatione pacis.|
|80.||In the trouble of peace||Gregory XV (1621–1623)||Alessandro Ludovisi||The lack of plausible explanations for this motto leads O'Brien to comment, "The prophet, up to 1590, did not deal in generalities."|
|Lilium et roſa.|
|81.||Lily and rose||Urban VIII (1623–1644)||Maffeo Barberini||This motto again may have been intended to suggest a heraldic device, but not one that matches Urban VIII's arms. Proponents of the prophecies have alternatively suggested that it is a reference to the bees that do occur on his arms, to the fleur-de-lis of his native Florence, or to his dealings in France (the lily) and England (the rose).|
|82.||Delight of the cross||Innocent X (1644–1655)||Giovanni Battista Pamphili||Proponents of the prophecies have attempted to link this motto to Innocent X by noting that he was raised to the pontificate around the time of the Feast of the Exaltation of the Cross.|
|83.||Guard of the mountains||Alexander VII (1655–1667)||Fabio Chigi||Proponents of the prophecies have attempted to link this motto to Alexander VII by noting that his papal arms include six hills, though this was not an uncommon device, and this explanation would not account for the "guard" portion of the motto.|
|84.||Star of the swans||Clement IX (1667–1669)||Giulio Rospigliosi||This again may have been intended to be taken as an allusion to heraldry; O'Brien notes that there is an Italian family with arms featuring a swan with stars, but it had no relation to Clement IX. Proponents of the prophecies have claimed he had a room called the "chamber of swans" during the conclave.|
|De flumine magno.|
|85.||From a great river||Clement X (1670–1676)||Emilio Altieri||Proponents of the prophecies have attempted to link this motto to Clement X by claiming that the Tiber overflowed its banks at his birth, or as an obscure reference to his family name.|
|86.||Insatiable beast||Innocent XI (1676–1689)||Benedetto Odescalchi||Proponents of the prophecies have attempted to link this motto to the lion on Innocent XI's arms.|
|87.||Glorious penitence||Alexander VIII (1689–1691)||Pietro Ottoboni||Proponents of the prophecies have attempted to link this motto to Alexander VIII by interpreting as a reference to the submission of the Gallican bishops. O'Brien notes, "There are glorious repentances during every pontificate."|
|Raſtrum in porta.|
|88.||Rake in the door||Innocent XII (1691–1700)||Antonio Pignatelli||Some sources discussing the prophecy give Innocent XII's family name as "Pignatelli del Rastello," which would provide a clear way for proponents to connect this motto to this pope (rastello or rastrello is Italian for rake). Others, however, give the pope's family name as simply "Pignatelli", and indicate that it is difficult to find a satisfactory explanation to associate the pope with the motto.|
|89.||Surrounded flowers||Clement XI (1700–1721)||Giovanni Francesco Albani||A medal of Clement XI was created with the motto, "Flores circumdati", drawn from his description in the prophecies, which were widely circulated at that time.|
|De bona religione.|
|90.||From good religion||Innocent XIII (1721–1724)||Michelangelo dei Conti||Proponents of the prophecies have attempted to link this motto to Innocent XIII by interpreting it as a reference to the fact several popes had come from his family.|
|Miles in bello.|
|91.||Soldier in War||Benedict XIII (1724–1730)||Pietro Francesco Orsini||Proponents of the prophecies have attempted to link this motto to particular wars that occurred during Benedict XIII's pontificate, or a figurative war against decadence in favour of austerity.|
|92.||Lofty column||Clement XII (1730–1740)||Lorenzo Corsini||This may have been intended by the author of the prophecies as a reference to a pope of the Colonna family; a similar motto was used to describe to Martin V, who was pope before the publication of the prophecies. Proponents of the prophecies have attempted to link this motto to Clement XII as an allusion to a statue erected in his memory or the use of two columns from the Pantheon of Agrippa in a chapel he built.|
|93.||Country animal||Benedict XIV (1740–1758)||Marcello Lambertini||This may have been intended as a reference to armorial bearings, but it does not match Benedict XIV's arms. Proponents of the prophecies have attempted to link this motto to this pope as a description of his "plodding ox" diligence.|
|94.||Rose of Umbria||Clement XIII (1758–1769)||Carlo Rezzonico||Proponents of the prophecies have attempted to link this motto to Clement XIII as a reference to his elevation to sainthood of several Franciscans, to which order the motto can refer.|
|95.||Swift bear (later misprinted as Cursus velox Swift Course or Visus velox Swift Glance)||Clement XIV (1769–1774)||Lorenzo Giovanni Vincenzo Antonio Ganganelli||Proponents of the prophecies have struggled to provide a satisfactory explanation of this motto; some authors claim without evidence that the Ganganelli arms featured a running bear, but this is dubious.|
|96.||Apostolic pilgrim||Pius VI (1775–1799)||Giovanni Angelico Braschi||Proponents of the prophecies have attempted to link this motto to Pius VI by suggesting it is a reference to his long reign.|
|97.||Rapacious eagle||Pius VII (1800–1823)||Barnaba Chiaramonti||Proponents of the prophecies have attempted to link this motto to Pius VII by suggesting it is a reference to the eagle on the arms of Napoleon, whose reign as Emperor of the French took place during Pius' pontificate.|
|Canis & coluber.|
|98.||Dog and adder||Leo XII (1823–1829)||Annibale Sermattei della Genga||Proponents of the prophecies have attempted to link this motto to Leo XII by suggesting the dog and snake are allusions to his qualities of vigilance and prudence, respectively.|
|99.||Religious man||Pius VIII (1829–1830)||Francesco Saverio Castiglioni||Proponents of the prophecies have attempted to link this motto to Pius VIII by suggesting it is a reference to his papal name, or the fact that he was not the first pope from his family.|
|De balneis Ethruriæ.|
|100.||From the baths of Tuscany||Gregory XVI (1831–1846)||Mauro, or Bartolomeo Alberto Cappellari||Proponents of the prophecies have attempted to link this motto to Gregory XVI by suggesting it is a reference to his membership in the Camaldolese Order, founded in the thirteenth century in Fonte Buono, called Balneum in Latin, in Etruria.|
|Crux de cruce.|
|101.||Cross from cross||Bl. Pius IX (1846–1878)||Giovanni Maria Mastai Ferretti||Proponents of the prophecies have attempted to link this motto to Pius IX by interpreting it as a reference to his difficulties ("crosses") with the House of Savoy, whose emblem is a cross. O'Brien notes, "A forger would be very disposed to chance some reference to a cross on account of its necessary connexion with all popes as well as the probability of its figuring, in some form or other, on the pope's arms."|
|Lumen in cœlo.|
|102.||Light in the sky||Leo XIII (1878–1903)||Gioacchino Pecci||Proponents of the prophecies have attempted to link this motto to Leo XIII by interpreting it as a reference to the star on his arms. O'Brien notes this coincidence would be much more remarkable had the prophecies referred to sydus (star), as they did when describing this same device on pre-publication Pope Innocent VII's arms.|
|103.||Burning fire||St. Pius X (1903–1914)||Giuseppe Sarto||Proponents of the prophecies have attempted to link this motto to Pius X by interpreting it as a reference to his zeal.|
|104.||Religion destroyed||Benedict XV (1914–1922)||Giacomo Della Chiesa||Proponents of the prophecies have attempted to link this motto to Benedict XV by interpreting it as a reference to World War I and the Russian Revolution, which occurred during his pontificate.|
|105.||Intrepid faith||Pius XI (1922–1939)||Achille Ratti||Proponents of the prophecies have attempted to link this motto to Pius XI by interpreting it as a reference to his faith and actions during the reign of Benito Mussolini.|
|106.||Angelic shepherd||Ven. Pius XII (1939–1958)||Eugenio Pacelli||Proponents of the prophecies have attempted to link this motto to Pius XII by interpreting it as a reference to his role during the holocaust.|
|Paſtor & nauta.|
|107.||Shepherd and sailor||Bl. John XXIII (1958–1963)||Angelo Giuseppe Roncalli||Proponents of the prophecies have attempted to link the "sailor" portion of this motto to John XXIII by interpreting it as a reference to his title Patriarch of Venice, a maritime city.|
|108.||Flower of flowers||Paul VI (1963–1978)||Giovanni Battista Enrico Antonio Maria Montini||Proponents of the prophecies have attempted to link this motto to Paul VI by interpreting it as a reference to the fleurs-de-lis on his arms.|
|De medietate lunæ.|
|109.||Of the half moon||John Paul I (1978)||Albino Luciani|
|De labore solis.|
|110.||From the labour of the sun / Of the eclipse of the sun||Bl. John Paul II (1978–2005)||Karol Wojtyła||Proponents of the prophecies find significance in the occurrence of solar eclipses (elsewhere in the world) on the dates of John Paul II's birth (18 May 1920) and funeral (8 April 2005). Other attempts to link the pope to the motto have been "more forced," included drawing a connection to Copernicus (who formulated a comprehensive heliocentric model of the solar system), as both were Polish and lived in Kraków for parts of their lives.|
|111.||Glory of the olive.||Benedict XVI (2005–2013)||Joseph Ratzinger||Proponents of the prophecies generally try to draw a connection between Benedict and the Olivetan order to explain this motto: Benedict's choice of papal name is after Saint Benedict of Nursia, founder of the Benedictine Order, of which the Olivetans are one branch. Other explanations make reference to him as being a pope dedicated to peace and reconciliations of which the olive branch is the symbol.|
|In perſecutione extrema S.R.E. ſedebit.|
|In the final persecution of the Holy Roman Church, there will sit.||In the Lignum Vitae, the line "In persecutione extrema S.R.E. sedebit." forms a separate sentence and paragraph of its own. While often read as part of the "Peter the Roman" prophecy, other interpreters view it as a separate, incomplete sentence explicitly referring to additional popes between "glory of the olive" and "Peter the Roman".|
|Petrus Romanus, qui paſcet oues in multis tribulationibus: quibus tranſactis ciuitas ſepticollis diruetur, & Iudex tremẽdus iudicabit populum ſuum. Finis.|
|112.||Peter the Roman, who will pasture his sheep in many tribulations, and when these things are finished, the city of seven hills [i.e. Rome] will be destroyed, and the dreadful judge will judge his people. The End.||Francis (2013–present)||Jorge Mario Bergoglio||Many analyses of the prophecy note that it is open to the interpretation that additional popes would come between the "glory of the olive" and Peter the Roman. Popular speculation by proponents of the prophecy attach this prediction to Benedict XVI's successor. Since Francis' election as Pope, proponents in internet forums have been striving to link him to the prophecy. Theories include a vague connection with Francis of Assisi, whose father was named Pietro (Peter).|
The Prophecies in art, literature, and culture
- Pope Patrick is a novel about the then Pope John Paul II's supposed successor, the fictional Pope Patrick I. The novel assumes that 'Petrus Romanus', the last Pope in St Malachy's list, is to be regarded as some kind of unreal supernatural being, and that consequently Pope Patrick will be the last real Pope.
- Glory of the Olive: A Novel of the Time of Tribulation (first published in 2002) is a novel featuring the fictional Pope Peter II. "Glory of the Olive" is the Malachite attribute of the successor to the then Pope John Paul II.
- The Roman: Peter II... The Last Pope? is a novel featuring the fictional Pope Peter II as successor to the then Pope Benedict XVI. The novel begins with the Malachite prophecy concerning Peter the Roman, the last Pope in the Malachite list.
- The Third Secret is a novel featuring the fictional Pope Peter II (originally Cardinal Valendrea), who is elected Pope after the death of the fictional Pope Clement XV.
- The Devil Will Come (novel) by Glenn Cooper is a novel which uses the Malachy Prophecy as a part of the storyline in the book which spans generations, leading to the "modern day" conclave to elect a new pope and the attempt to destroy the Catholic faith by an enemy of the church.
List of fictional Popes Peter II
(Note: The final Pope in the Malachite list is called Peter the Roman)
- Pope Peter II – War, Progress, and the End of History: Three Conversations, Including a Short Story of the Anti-Christ by Vladimir S. Solovyov
- Pope Peter II – Petrus Secundus by Harold J. Frysne
- Pope Peter II – Peter the Second by Bruce Marshall (Third part of a Trilogy)
- Pope Peter II – The Accidental Pope by Raymond Flynn and Robin Moore
- Pope Peter II – The Final Restoration by John Cantwell Kiley
- Pope Peter II – The Reckoning by Thomas F. Monteleone
- Pope Peter II – The Third Secret: A Novel Of Suspense by Steve Berry
- Pope Peter II – Left Behind novels by Tim LaHaye and Jerry B. Jenkins
- Pope Peter II – Glory of the Olive: A Novel of the Time of Tribulation by Susan Claire Potts
- Pope Peter II – En el mar de la duda by Pedro De Illanez
- Pope Peter II – The Roman: Peter II... The Last Pope? by George R. Araujo-Matiz
- Pope Peter II – Fumata Bianca by Giuseppe Magnarapa
- Pope Peter II – I giorni della tempesta by Antonio Socci
- Pope Peter II (Antipope)– L'évangile selon Satan by Patrick Graham
- Bartholomew Holzhauser
- Bible code
- Legends surrounding the papacy
- List of popes
- The Prophesying Nun of Dresden
- Three Secrets of Fátima
- Vaticinia de Summis Pontificibus
- Vaticinia Nostradami
- Sieczkowski 2013.
- Boyle 2013.
- O'Brien 1880, pp. 16 & 25.
- Menestrier 1694, pp. 343-344.
- Catholic Encyclopedia 1913, "Prophecy".
- O'Brien 1880, p. 110.
- de Vallemont 1708, p. 87.
- Feijóo y Montenegro 1724-1739, p. 129.
- O'Brien 1880, p. 14.
- O'Brien 1880, p. 85.
- Feijóo y Montenegro 1724-1739, p. 134.
- Allan 2009, pp. 58-9.
- "Petrus Romanus Prophecy; Will The Next Pope Lead To The Apocalypse?", from International Business Times
- See, e.g. Bander 1969, p. 96.
- O'Brien 1880, p. 82.
- René Thibaut S.J.: La mystérieuse prophétie des Papes. Namur-Paris, 1951, p. 10.
- Hildebrand Troll: Die Papstweissagung des heiligen Malachias. Ein Beitrag zur Lösung ihres Geheimnisses. EOS-Verlag, St. Ottilien 2002, ISBN 3-8306-7099-0.
- O'Brien 1880, p. 47.
- O'Brien 1880, p. 28.
- O'Brien 1880, p. 28; Bander 1969, p. 19.
- Dizionario Biografico degli Italiani 2007, "Lucio II, papa".
- O'Brien 1880, p. 29; Bander 1969, p. 19.
- Dizionario Biografico degli Italiani 2007, "Eugenio III, papa".
- Michael Horn, Studien zur Geschichte Papst Eugens III.(1145-1153), Peter Lang Verlag 1992, pp. 28-33.
- Enciclopedia dei papi Treccani
- Hüls, Rudolf: Kardinäle, Klerus und Kirchen Roms: 1049–1130. Bibliothek des Deutschen Historischen Instituts in Rom. Max Niemeyer Verlag. Tübingen 1977,p. 201. ISBN 978-3-484-80071-7
- O'Brien 1880, p. 31.; Bander 1969, p. 23.
- O'Brien 1880, p. 31; Bander 1969, p. 25.
- O'Brien 1880, p. 33; Bander 1969, p. 26.
- Johannes Matthias Brixius, Die Mitglieder des Kardinalkollegiums von 1130-1181. Berlin : R. Trenkel, 1912, p. 68-69, no. 1
- O'Brien 1880, p. 34; Bander 1969, p. 24.
- O'Brien 1880, p. 36; Bander 1969, p. 24.
- O'Brien 1880, p. 36; Bander 1969, p. 28.
- O'Brien 1880, p. 37; Bander 1969, p. 28.
- A non-standard verb form, replacing classical exibit.
- O'Brien 1880, p. 37; Bander 1969, p. 29.
- Bander 1969, p. 30.
- O'Brien 1880, p. 38; Bander 1969, p. 30.
- O'Brien 1880, p. 39; Bander 1969, p. 32.
- O'Brien 1880, p. 40; Bander 1969, p. 33.
- O'Brien 1880, p. 40; Bander 1969, p. 34.
- O'Brien 1880, p. 41; Bander 1969, p. 35.
- O'Brien 1880, p. 42; Bander 1969, p. 35.
- O'Brien 1880, p. 42; Bander 1969, p. 36.
- O'Brien 1880, p. 43; Bander 1969, p. 36.
- O'Brien 1880, p. 43; Bander 1969, p. 37.
- Bander 1969, p. 38.
- O'Brien 1880, p. 44.
- O'Brien 1880, p. 44; Bander 1969, p. 39.
- Properly Asculanus, but that ruins the pun.
- O'Brien 1880, p. 45; Bander 1969, p. 41.
- O'Brien 1880, p. 46; Bander 1969, p. 42.
- O'Brien 1880, p. 47; Bander 1969, p. 43.
- O'Brien 1880, p. 48; Bander 1969, p. 44.
- O'Brien 1880, p. 48; Bander 1969, p. 45.
- O'Brien 1880, p. 49; Bander 1969, p. 45.
- O'Brien 1880, p. 49; Bander 1969, p. 46.
- O'Brien 1880, p. 49; Bander 1969, p. 47.
- Bander 1969, p. 47.
- O'Brien 1880, p. 50.
- O'Brien 1880, p. 50; Bander 1969, p. 48.
- O'Brien 1880, p. 51; Bander 1969, p. 50.
- O'Brien 1880, p. 52; Bander 1969, p. 51.
- O'Brien 1880, p. 53; Bander 1969, p. 48.
- O'Brien 1880, p. 53; Bander 1969, p. 49.
- O'Brien 1880, p. 54; Bander 1969, p. 50.
- O'Brien 1880, p. 54; Bander 1969, p. 52.
- O'Brien 1880, p. 55; Bander 1969, p. 53.
- O'Brien 1880, p. 55; Bander 1969, p. 54.
- O'Brien 1880, p. 56; Bander 1969, p. 56.
- O'Brien 1880, p. 56; Bander 1969, p. 57.
- O'Brien 1880, p. 57; Bander 1969, p. 58.
- O'Brien 1880, p. 57; Bander 1969, p. 59.
- O'Brien 1880, p. 58; Bander 1969, p. 60.
- O'Brien 1880, p. 58; Bander 1969, p. 61.
- O'Brien 1880, p. 58; Bander 1969, p. 62.
- Pileus here is not usually translated as "cap", but as if derived from pila "ball" or Late Latin pilula "little ball, pill".
- O'Brien 1880, p. 59; Bander 1969, p. 62.
- O'Brien 1880, p. 59; Bander 1969, p. 63.
- O'Brien 1880, p. 60; Bander 1969, p. 64.
- O'Brien 1880, p. 60; Bander 1969, p. 65.
- O'Brien 1880, p. 61; Bander 1969, p. 66.
- O'Brien 1880, p. 61; Bander 1969, p. 67.
- O'Brien 1880, p. 61; Bander 1969, p. 68.
- O'Brien 1880, p. 62; Bander 1969, p. 68.
- O'Brien 1880, p. 62; Bander 1969, p. 70.
- O'Brien 1880, p. 63; Bander 1969, p.70.
- O'Brien 1880, p. 64; Bander 1969, p.71.
- O'Brien 1880, p. 64; Bander 1969, p.72.
- O'Brien 1880, p. 65; Bander 1969, p.72.
- O'Brien 1880, p. 65.
- O'Brien 1880, p. 66.
- O'Brien 1880, p. 66; Bander 1969, p.75.
- O'Brien 1880, p. 67; Bander 1969, p.75.
- O'Brien 1880, p. 67; Bander 1969, p.76.
- O'Brien 1880, p. 69.
- O'Brien 1880, p. 69; Bander 1969, p. 77.
- O'Brien 1880, p. 70; Bander 1969, p. 78.
- Bander 1969, p. 79.
- O'Brien 1880, p. 70
- Rastellus, a diminutive of rastrum, can also refer to a metallic grid used to close the door of a town during night, cataracta in portis urbium according to Du Cange et al, Glossarium mediae et infimae latinitatis, ad vocem.
- See, e.g., de Vallemont 1708, p. 123, and Cucherat 1873, p. 206 (citing de Vallemont).
- O'Brien 1880, p. 70; Bander 1969, p. 79.
- O'Brien 1880, p. 71; Bander 1969, p. 79.
- O'Brien 1880, p. 71; Bander 1969, p. 80.
- O'Brien 1880, p. 72; Bander 1969, p. 80.
- O'Brien 1880, p. 72; Bander 1969, p. 81.
- O'Brien 1880, p. 73; Bander 1969, p. 83.
- O'Brien 1880, p. 74; Bander 1969, p. 83.
- O'Brien 1880, p. 74; Bander 1969, p. 84.
- The symbol like a raised 9 is a scribal abbreviation for the Latin suffix us.
- O'Brien 1880, p. 75; Bander 1969, p. 85.
- O'Brien 1880, p. 75; Bander 1969, p. 86.
- O'Brien 1880, p. 77; Bander 1969, p. 87.
- O'Brien 1880, p. 76; Bander 1969, p. 87.
- O'Brien 1880, p. 78; Bander 1969, p. 88.
- O'Brien 1880, p. 79; Bander 1969, p. 89.
- Bander 1969, p. 90.
- Bander 1969, p. 91; Allan 2009, pp. 58-9.
- Bander 1969, p. 91.
- Bander 1969, p. 92.
- Bander 1969, p. 93.
- Bander 1969, p. 94; Allan 2009, pp. 58-9.
- O'Brien 1880, p. 81.
- Bander 1969, p. 94.
- Bander 1969, p. 95.
- Gloria Olivae as a Peace Symbol, Does Pope Benedict XVI's resignation signal the 'end times?'
- In several later printings of the prophecies, the word ſuum was dropped, leading to the translation "the people" instead of "his people". See, e.g., O'Brien 1880, p. 83.
- "Forums strive to connect new Pope to Antichrist prophecy", from The Fraser Coast Chronicle
- De Rosa, Peter (1997). "''Pope Patrick''". Doubleday. ISBN 978-0385485487. Retrieved 2013-02-17.
- Potts, Susan Claire (2002). "Glory of the Olive: A Novel of the Time of Tribulation". iUniverse. ISBN 978-0595223220. Retrieved 2013-02-17.
- Araujo-Matiz, George R. (2007-02-21). "The Roman: Peter II... The Last Pope?". BookSurge Publishing. ISBN 978-1419651403. Retrieved 2013-02-17.
- Berry, Steve. "The Third Secret".
- Devil will come Glenn Cooper, Harpercollins.ca
- Allan, Tony (2009). Prophecies: 4,000 years of prophets, visionaries and predictions. London: Duncan Baird. ISBN 1780283407. Retrieved 12 February 2013.
- Boyle, Alan (12 February 2013). "Why the buzz over St. Malachy's 'last pope' prophecy outdoes 2012 hype". NBC News. Retrieved 17 February 2013.
- Bander, Peter (1969). The Prophecies of St. Malachy. Buckinghamshire, England: Colin Symthe Ltd.
- Cucherat, François (1873). La prophétie de la succession des papes depuis le XIIe siècle jusqu'a la fin du monde (in French). Grenoble: E. Dardelet.
- de Vallemont, Pierre Le Lorrain (1708). Les élemens de l'histoire ou ce qu'il faut savoir (in French) 3. Paris: Chez Rigaud, Directeur de l'Imprimerie Royale.
- "Eugenio III, papa". Dizionario Biografico degli Italiani (in Italian) 43. Istituto dell'Enciclopedia Italiana. 1993. Retrieved 19 February 2013.
- Feijóo y Montenegro, Benito Jerónimo (1724-1739). Teatro crítico universal (in Spanish). p. 129.
- "Lucio II, papa". Dizionario Biografico degli Italiani (in Italian) 66. Istituto dell'Enciclopedia Italiana. 2007. Retrieved 19 February 2013.
- Menestrier, Claude-François (1694). La philosophie des images énigmatiques, où il est traité des énigmes hiéroglyphiques, oracles, prophéties, sorts etc (in French). University of Lausanne.
- O'Brien, M. J. (1880). An historical and critical account of the so-called Prophecy of St. Malachy, regarding the succession of the popes. Dublin: M.H. Gill & Son.
- "Prophecy". Catholic Encyclopedia. New Advent. 1913. Retrieved 12 February 2013.
- Sieczkowski, Cavan (14 February 2013). "St. Malachy Last Pope Prophecy: What Theologians Think About 12th-Century Prediction". Huffington Post Canada. Retrieved 17 February 2013.
- Original 1595 text of the Prophecies (Arnold Wion, Lignum Vitae, Lib. ii, pp. 307–311) |
United States congressional apportionment
United States congressional apportionment is the process by which seats in the United States House of Representatives are redistributed amongst the 50 states following each constitutionally mandated decennial census. Each state is apportioned a number of seats which approximately corresponds to its share of the aggregate population of the 50 states. However, every state is constitutionally guaranteed at least one seat.
The number of seats in the House of Representatives is currently set to 435, and has been since 1913, except for a temporary increase to 437 after the admissions of Alaska and Hawaii. Though the actual reapportionment will normally occur in respect of a decennial census, the law that governs the total number of representatives and the method of apportionment to be carried into force at that time can be created prior to the census.
The decennial apportionment also determines the size of each state's representation in the U.S. Electoral College—that is, any state's number of electors equals the size of its total congressional delegation (i.e., House seat(s) plus Senate seats).
Federal law requires the Clerk of the House of Representatives to notify each state government of its entitled number of seats no later than January 25 of the year immediately following the census. After seats have been reapportioned, each state determines the boundaries of congressional districts—geographical areas within the state of approximately equal population—in a process called redistricting. Any citizen of the State can challenge the constitutionality of the redistricting in their US district court.
Because the deadline for the House Clerk to report the results does not occur until the following January, and the states need sufficient time to perform the redistricting, the decennial census does not affect the elections that are held during that same year. For example, the electoral college apportionment during 2000 presidential election was still based on the 1990 census results. Likewise, the congressional districts and the electoral college during the 2020 general elections will still be based on the 2010 census.
Constitutional text
Representatives and direct Taxes shall be apportioned among the several States which may be included within this Union, according to their respective Numbers, which shall be determined by adding to the whole Number of free Persons, including those bound to Service for a Term of Years, and excluding Indians not taxed, three fifths of all other Persons.
Representatives shall be apportioned among the several States according to their respective numbers, counting the whole number of persons in each State, excluding Indians not taxed.
Article I additionally provides that:
The Number of Representatives shall not exceed one for every thirty Thousand, but each State shall have at Least one Representative...
House size
The size of the U.S. House of Representatives refers to total number of congressional districts (or seats) into which the land area of the United States proper has been divided. The number of voting representatives is currently set at 435. There are an additional five delegates to the House of Representatives. They represent the District of Columbia and the territories of American Samoa, Guam, the Northern Mariana Islands, which first elected a representative in 2008, and the U.S. Virgin Islands. Puerto Rico also elects a resident commissioner every four years.
Controversy and history
During the period that the current U.S. Constitution has been in effect, the number of citizens per congressional district has risen from an average of 33,000 in 1790 to almost 700,000 as of 2008[update]. Prior to the 20th century, the number of representatives increased every decade as more states joined the union, and the population increased. In 1911, Public Law 62-5 raised the membership of the U.S. House to 433 with a provision to add one permanent seat each upon the admissions of Arizona and New Mexico as states. As provided, membership increased to 435 in 1912.
But in 1921, Congress failed to reapportion the House membership as required by the United States Constitution. This failure to reapportion may have been politically motivated, as the newly elected Republican majority may have feared the effect such a reapportionment would have on their future electoral prospects. Then in 1929 Congress (Republican control of both houses of congress and the presidency) passed the Reapportionment Act of 1929 which capped the size of the House at 435 (the then current number), but allowed temporary increases upon the admission of new states which were to be reverted upon the implementation of the immediate subsequent census.
In truth, the rules prohibiting legislative entrenchment would allow any subsequent legislature (after 1929) to increase or decrease the membership of the House of Representatives if such legislature so desired.
The current size of 435 seats means one member represents on average about 709,760 people; but exact representation per member varies by state. Three states – Wyoming, Vermont, and North Dakota – have populations smaller than the average for a single district.
The "ideal" number of members has been a contentious issue since the country's founding. George Washington agreed that the original representation proposed during the Constitutional Convention (one representative for every 40,000) was inadequate and supported an alteration to reduce that number to 30,000. This was the only time that Washington pronounced an opinion on any of the actual issues debated during the entire convention. At the same time, the researchers with RangeVoting.org contend the optimal legislature size should be somewhere between the cube root and the square root of the constituent population, which would yield a number between 700 and 17,000 for the current U.S. population.
In Federalist No. 55, James Madison argued that the size of the House of Representatives has to balance the ability of the body to legislate with the need for legislators to have a relationship close enough to the people to understand their local circumstances, that such representatives' social class be low enough to sympathize with the feelings of the mass of the people, and that their power be diluted enough to limit their abuse of the public trust and interests.
"... first, that so small a number of representatives will be an unsafe depositary of the public interests; secondly, that they will not possess a proper knowledge of the local circumstances of their numerous constituents; thirdly, that they will be taken from that class of citizens which will sympathize least with the feelings of the mass of the people, and be most likely to aim at a permanent elevation of the few on the depression of the many;..."
Madison also addressed Anti-Federalist claims that the representation would be inadequate, arguing that the major inadequacies are of minimal inconvenience since these will be cured rather quickly by virtue of decennial reapportionment. He noted, however,
"I take for granted here what I shall, in answering the fourth objection, hereinafter show, that the number of representatives will be augmented from time to time in the manner provided by the Constitution. On a contrary supposition, I should admit the objection to have very great weight indeed."
Madison argued against the assumption that more is better:
"Sixty or seventy men may be more properly trusted with a given degree of power than six or seven. But it does not follow that six or seven hundred would be proportionally a better depositary. And if we carry on the supposition to six or seven thousand, the whole reasoning ought to be reversed. ... In all very numerous assemblies, of whatever character composed, passion never fails to wrest the scepter from reason."
Clemons v. Department of Commerce
A recent lawsuit, Clemons v. Department of Commerce, sought a court order for Congress to increase the size of the House's voting membership and then reapportion the seats in accordance with the population figures of the 2010 Census. The intent of the plaintiff was to rectify the disparity of congressional district population sizes among the states that result from the present method of apportionment. Upon reaching the U.S. Supreme Court in December 2010, the holdings of the lower district and appellate courts were vacated and the case remanded to the U.S. District Court from which the case originated with instructions that the district court dismiss the case for lack of jurisdiction.
Proposed expansion
The proposed Wyoming Rule calls for expanding the House until the standard Representative-to-population ratio equals that of the smallest entitled unit (currently the state of Wyoming). This proposal is primarily designed to address the fact that some House districts are currently nearly twice the size of others; for instance, there are nearly 1 million residents in Montana's single district, compared to about 570,000 in Wyoming's. See List of U.S. states by population.
In 2007, during the 110th Congress, Representative Tom Davis introduced a bill in the House of Representatives that would add two seats to the House, one for Utah and one for the District of Columbia. It was passed by the House, but was tripped up by procedural hurdles in Senate and withdrawn from consideration. An identical bill was reintroduced during the 111th Congress. In February 2009 the Senate adopted the measure 61-37. In April 2010, however, House leaders decided to shelve the proposal.
Apportionment methods
Apart from the requirement that the number of delegates for each state is at least one, a state's number of representatives is in principle proportional to population, thus assuring reasonably consistent representation to the people regardless of the state boundaries and populations. No method of calculating a fair distribution of voting power across the various states was known until recently and five distinct apportionment methods have been used since the adoption of the Constitution, none of them producing fully proportional distribution of power among the states. Some of these methods were even susceptible to the apportionment paradox. In 2008, however, a so-called One-Person-One-Vote model was introduced by J. Toplak in Temple Law Review, which distributes the states' power in the House of Representatives exactly 'according to their Numbers'. Under this system, however, members of the House of Representatives do not have equal voting power. The method would be constitutional since the U.S. Constitution does not require Congressmen to have equal voting powers but does require the voters to have votes of equal weight.
The Method of Equal Proportions
The apportionment methodology currently used is the method of equal proportions, so called because it guarantees that no additional transfer of a seat (from one state to another) will reduce the ratio between the numbers of persons per representative in any two states. According to NationalAtlas.gov, the method of equal proportions minimizes the percentage differences in the populations of the congressional districts.
In this method, as a first step, each of the 50 states is given its one guaranteed seat in the House of Representatives, leaving 385 seats to assign.
The remaining seats are allocated one at a time, to the state with the highest priority number. Thus, the 51st seat would go to the most populous state (currently California). The priority number is determined by a formula that is mathematically computed to be the ratio of the state population to the geometric mean of the number of seats it currently holds in the assignment process, n (initially 1), and the number of seats it would hold if the seat were assigned to it, n+1.
The formula for determining the priority of a state to be apportioned the next available seat defined by the method of equal proportions is
where P is the population of the state, and n is the number of seats it currently holds before the possible allocation of the next seat. An equivalent, recursive definition is
where n is still the number of seats the state has before allocation of the next, and for n = 1, the initial A1 is explicitly defined as
Consider the reapportionment following the 2010 U.S. Census: beginning with all states initially being allocated one seat, the largest value of A1 corresponds to the largest state, California, which is allocated seat 51. After being allocated its 2nd seat, its priority value decreases to its A2 value, which is reordered to a position back in line. The 52nd seat goes to Texas, the 2nd largest state, because its A1 priority value is larger than the An of any other state. However, the 53rd seat goes back to California because its A2 priority value is larger than the An of any other state. The 54th seat goes to New York because its A1 priority value is larger than the An of any other state at this point. This process continues until all remaining seats are assigned. Each time a state is assigned a seat, n is incremented by 1, causing its priority value to be reduced and reordered among the states, whereupon another state normally rises to the top of the list.
The Census 2010 Ranking of Priority Values shows the order in which seats 51–435 were apportioned after the 2010 Census, with additional listings for the next five priorities. Minnesota was allocated the final (435th) seat. North Carolina, which gained the 435th seat in the 2000 census, missed its 14th seat by 15,754 residents as the 436th seat to be allocated.
Past apportionments
Note: The first apportionment was established by the Constitution based on population estimates made by the Philadelphia Convention, and was not based on any census or enumeration.
|14th, 1920||Congress failed to pass any reapportionment act. Distribution of seats from 1913 remained in effect, despite population shifts.|
Changes following the 2010 census
|Gain four||Gain two||Gain one||Lose one||Lose two|
Past increases
|March 4, 1789||59||n/a||Const. Art. I,
§ 2, cl. 3
|Seats apportioned by the Constitution|
|November 21, 1789||64||5||North Carolina ratified the Constitution with the seats apportioned by the Constitution|
|May 29, 1790||65||1||Rhode Island ratified the Constitution with the seats apportioned by the Constitution|
|March 4, 1791||67||2||1 Stat. 191||Vermont admitted|
|June 1, 1792||69||2||Kentucky admitted|
|March 4, 1793||105||36||1 Stat. 253||Apportionment of the First Census|
|June 1, 1796||106||1||1 Stat. 492||Tennessee admitted|
|March 1, 1803||107||1||2 Stat. 175||Ohio admitted.|
|March 4, 1803||142||35||2 Stat. 128||Apportionment of the Second Census.|
|April 30, 1812||143||1||2 Stat. 703||Louisiana admitted.|
|March 4, 1813||182||39||2 Stat. 669||Apportionment of the Third Census.|
|December 11, 1816||183||1||3 Stat. 290||Indiana admitted.|
|December 10, 1817||184||1||3 Stat. 349||Mississippi admitted.|
|December 3, 1818||185||1||3 Stat. 430||Illinois admitted.|
|December 14, 1819||186||1||3 Stat. 492||Alabama admitted.|
|March 15, 1820||3 Stat. 555||Maine admitted, 7 seats transferred from Massachusetts|
|August 10, 1821||187||1||3 Stat. 545||Missouri admitted|
|March 4, 1823||213||26||3 Stat. 651||Apportionment of the Fourth Census|
|March 4, 1833||240||27||4 Stat. 516||Apportionment of the Fifth Census|
|June 15, 1836||241||1||5 Stat. 49||Arkansas admitted|
|January 26, 1837||242||1||5 Stat. 50||Michigan admitted|
|March 4, 1843||223||19||5 Stat. 491||Apportionment of the Sixth Census, the only time the size of the House was reduced, except for the minor readjustment in 1963.|
|March 3, 1845||224||1||5 Stat. 743||Florida admitted.|
|December 29, 1845||226||2||5 Stat. 798||Texas annexed and admitted.|
|December 28, 1846||228||2||5 Stat. 743
9 Stat. 52
|May 29, 1848||230||2||9 Stat. 58
9 Stat. 235
|March 4, 1849||231||1||9 Stat. 235||Wisconsin given another seat.|
|September 9, 1850||233||2||9 Stat. 452||California admitted.|
|March 4, 1853||9 Stat. 432||Apportionment of the Seventh Census.|
|234||1||10 Stat. 25||Additional seat apportioned to California|
The 1850 Apportionment bill provided a method to be used in future reapportionments, as well as establishing the then-current 233 as the number of seats to be apportioned after future censuses. Due to census returns being incomplete in California, an additional act provided that California retain the same representation it had when admitted, until a new census could be taken. California would otherwise have lost one seat, and so the total number of seats was increased by one to 234.
|May 11, 1858||236||2||11 Stat. 166||Minnesota admitted.|
|February 14, 1859||237||1||11 Stat. 383||Oregon admitted.|
|January 29, 1861||238||1||11 Stat. 269||Kansas admitted|
|June 2, 1862||239||1||12 Stat. 411||California apportioned an extra seat|
|March 4, 1863||233||6||9 Stat. 432||Apportionment of the Eighth Census, in accordance with the 1850 act, which provided for an apportionment of 233 seats|
|241||8||12 Stat. 353||Supplemental apportionment of 8 seats, for an overall increase of 2 seats in the 38th Congress|
|June 20, 1863||12 Stat. 633||West Virginia admitted, three seats transferred from Virginia|
|October 31, 1864||242||1||13 Stat. 32||Nevada admitted|
|March 1, 1867||243||1||14 Stat. 391||Nebraska admitted|
|March 4, 1873||283||40||17 Stat. 28||Apportionment of the Ninth Census, replacing the 1850 act|
|292||9||17 Stat. 192||Supplemental apportionment added one seat each for nine states|
|August 1, 1876||293||1||13 Stat. 34||Colorado admitted|
|March 4, 1883||325||32||47th Congress, ch. 20, 22 Stat. 5||Apportionment of the Tenth Census.|
|November 2, 1889||328||3||North Dakota and South Dakota admitted. One seat goes to the former, two to the latter.|
|November 8, 1889||329||1||Montana admitted.|
|November 11, 1889||330||1||Washington admitted.|
|July 3, 1890||331||1||Idaho admitted.|
|July 10, 1890||332||1||Wyoming admitted.|
|March 4, 1893||356||24||Apportionment Act of 1891||Apportionment of the Eleventh Census.|
|January 4, 1896||357||1||Utah admitted.|
|March 4, 1903||386||29||Apportionment Act of 1901, 31 Stat. 733||Apportionment of the Twelfth Census|
|November 16, 1907||391||5||Oklahoma Enabling Act||Oklahoma admitted|
|January 6, 1912||393||2||Apportionment Act of 1911, Sec. 2; New Mexico Enabling Act||New Mexico admitted|
|February 14, 1912||394||1||Apportionment Act of 1911, Sec. 2||Arizona admitted|
|March 4, 1913||435||41||Apportionment Act of 1911, Pub.L. 62–5, 37 Stat. 13||Apportionment of the Thirteenth Census|
|January 3, 1959||436||1||Alaska Statehood Act||Alaska admitted|
|August 21, 1959||437||1||Hawaii Admission Act||Hawaii admitted|
|January 3, 1963||435||2||Reapportionment Act of 1929, ch. 28, 46 Stat. 26, 2 U.S.C. § 2a||Apportionment of the Eighteenth Census
Note: The Reapportionment Act of 1929 stated that the "then existing number of Representatives" would be apportioned after each census, which would have dictated an apportionment of 437 seats, but the Alaska Statehood Act and Hawaii Admission Act explicitly stated that the new seats were temporary increases
See also
- Apportionment paradox
- Article the First
- List of U.S. states by population
- List of U.S. states by historical population (tables of state populations since 1790)
- United States Congress
- Delegate counts in italics represent temporary counts assigned by Congress until the next decennial census or by the U.S. Constitution in 1789 until the first U.S. Census.
- Elections held in the year of a census use the apportionment determined by the previous census.
- The populations of Washington, D.C. and federal territories are not included in this figure.
- U.S. Constitution, Article II, Section 1, Clause 2.
- 2 U.S.C. § 2c
- U.S. Const., art. I, § 2, cl. 3.
- U.S. Const., amend. XIV, § 2.
- Bush signs federalization bill, Agnes E. Donato, Saipan Tribune, May 10, 2008.
- "Fair Representation, Meeting The Ideal of One Man One vote" - Michael Balinski and H. Peyton Young -- Page 51
- George Will Called Me An Idiot, Jonah Golderg, National Review, January 15, 2001.
- Madison's notes on the Constitutional Convention - Tuesday September 17, 1787
- The Federalist #55
- The Federalist #55
- House of Representatives? Hardly., Alcee Hastings, May 21, 2001.
- Marimow, Ann E.; Pershing, Ben (April 21, 2010). "Congressional leaders shelve D.C. voting rights bill". The Washington Post.
- Toplak, Jurij (2008). "Equal Voting Weight of All: Finally 'One Person, One Vote' from Hawaii to Maine?". Temple Law Review (Temple University) 81 (1): 123–176.
- "The History of Apportionment in America". American Mathematical Society. Retrieved 2009-02-15.
- "2 USC §2a". Cornell University Law School, Legal Information Institute. Retrieved 2008-05-13.
- "Computing Apportionment". U.S. Census Bureau. Retrieved 2009-02-14.
- Edward V Huntington (1921). "The Mathematical Theory of the Apportionment of Representatives". Proceedings of the National Academy of Sciences U.S.A. 7 (4): 123–7. PMC 1084767. PMID 16576591.
- "Congressional Apportionment". NationalAtlas.gov. Retrieved 2009-02-14.
- "PRIORITY VALUES FOR 2010 CENSUS". U.S. Bureau of the Census. Retrieved 2012-06-07.
- "Census 2000 Ranking of Priority Values". U.S. Bureau of the Census. 2001-02-21. Retrieved 2008-05-13.
- "APPORTIONMENT POPULATION AND NUMBER OF REPRESENTATIVES, BY STATE: 2010 CENSUS". US Census. 2010-12-21. Retrieved 2013-02-23.
- The Size of the U. S. House of Representatives and its Constituent State Delegations, thirty-thousand.org.
- Both acts included the phrasing That such temporary increase in the membership shall not operate to either increase or decrease the permanent membership of the House of Representatives as prescribed in the Act of August 8, 1911 (37 Stat. 13) nor shall such temporary increase affect the basis of apportionment established by the Act of November 15, 1941 (55 Stat. 761; 2 U. S. C., sec. 2a), for the Eighty-third Congress and each Congress thereafter.
- Balinski, Michael L.; Young, H. Peyton (1982). Fair Representation: Meeting the Ideal of One Man, One Vote. New Haven, CT: Yale University Press. ISBN 0-8157-0090-3.
- Foster, Robert (1895). Commentaries on the Constitution of the United States: Historical and Judicial 1. Boston: The Boston Book Co. pp. 329–446.
- Hamilton, Alexander; Madison, James; Jay, John (1831). The Federalist. Hallowell: Glazier, Masters & Co. ISBN 0-8239-5735-7.
- Edelman, Paul H. (2006). "Getting the Math Right: Why California Has Too Many Seats in the House of Representatives". Vanderbilt Law Review (Nashville: Vanderbilt University) 102 (March): 297.
- Charles Kromkowski and John A Kromkowski; Kromkowski, John A (1991). "Why 435? A Question of Political Arithmetic" (PDF). Polity 24 (Fall 1991): 129–145. doi:10.2307/3234988. JSTOR 3234988. Retrieved 2009-06-22. More than one of
- Agnew, Robert A. (2008). "Optimal Congressional Apportionment". American Mathematical Monthly (Mathematical Association of America) 115 (April): 297–303.
Further reading
- Stinebrickner-Kauffman, Taren (2004). "Counting Matters: Prison Inmates, Population Bases, and "One Person, One Vote"". Virginia Journal of Social Policy & the Law (Chicago: Virginia Journal of Social Policy & the Law) 11 (Winter): 229. |
OSMap: Hadrian's Wall, OL43, LR87.
Type: Vexillation Fortress, Fort, Town, Bridge, Mausoleum.
Dere Street: N (2½) to Onnvm (Halton Chesters, Northumberland)|
Stanegate: W (6¾) to Cilvrnvm (Chesters, Northumberland)
Stanegate: ESE (15) to Washing Wells (Whickham, Tyne & Wear)
Dere Street: SE (9) to Vindomora (Ebchester, Durham)
The most striking thing about the Corbridge site, aside from the massive remains themselves, is the undulating ground caused by subsidence following the Roman departure, giving the appearance that the ruined town was originally built upon the green and gentle waves of a turf sea. It also makes the layout of the town very difficult to see from ground level unless one looks directly along the line of the Stanegate itself, or one of its side-streets.
Strategically placed beside the lowest fordable point of the Tyne, close downstream of the confluence of its North and South streams. Dere Street crossed the river here on its journey north from Eburacum (York) and continued northwards across the river into Barbaricum.
"Further south are the Otalini, among whom are the following towns: Coria 20*10 59°00, Alauna 23*00 58°40, Bremenium 21*00 58°45."
There is an interesting passage in the early second century geographical treatise by Claudius Ptolemaeus (see above), in which he assigns a town named Coria to the Otalini or Otadini tribe, along with other towns Alauna and Bremenium. These settlements have all been identified with Roman sites in Northumberland, at Corbridge, Learchild and High Rochester respectively. The tribal name has now been equated with the Votadini, whose territories lay mostly within the county of Northumberland in north-east England, also in the Borders region of south-east Scotland.
There is further mention of the Roman name for Corbridge in the Antonine Itinerary of the late-second century. The first route listed for Britain in this document is entitled "From the Wall at the limits [of the empire], to Praetorium", which lists the road-stations along the route from Hadrian's Wall to Bridlington on the north-east coast of England. The second entry in Iter I is named Corstopitum, and is listed 20 miles from Bremenium (High Rochester, Northumberland) and 9 miles from Vindomora (Ebchester, Durham).
There is further mention of Corbridge in the Ravenna Cosmography of the fourth/fifth century, wherein entry #142 appears as Corie Lopocarium, the first part of which has been equated with the Coria of Ptolemy (vide supra), but the suffix Lopocarium remains a mystery, and indeed, may even be the name of another town, as yet unidentified. The Corie entry appears between those of Concangis (Chester-le-Street, Durham) and Segedunum (Wallsend, Tyne & Wear).
"Corbridge Northum. Corebricg c.1050. 'Bridge near Corchester'. OE brycg 'bridge' with a shortened form of the old Celtic name of Corchester (Corstopitum) which is near here." (Mills, 1998)
The names Coria from Ptolemy and Corie from the RC, may be derived from the same Celtic roots as the Gaelic word Coire 'a round hollow in a mountainside', and the Welsh word Cwm 'valley, dale'; both words adequately describe the location of the Corbridge station. The Antonine name Corstopitum, is possibly a Romanisation of the original Celtic name suffixed by the Latin word strepitum 'loud noise, resounding', the Roman-British name therefore meaning something along the lines of 'The Valley of the Resounding Noise', a name which undoubtedly reflects its use as a busy legionary garrison post close to the troublesome Scottish border region. It may be significant that the entry identified with Corbridge in the Ravenna Cosmography, seems to show that by the seventh century the settlement had reverted back to its original Celtic name, without the Latin suffix.
|I O M PRO SALVTE VEXILLATIO NVM LEG XX V V ET VI VIC MILITES AGENT IN P...|
|"To Jupiter Best and Greatest, for the well-being and harmony of the Vexillation drawn from the Valiant and Victorious Twentieth Legion and the Victorious Sixth Legion, the soldiers negotiated the placing [of this]."|
|(RIB 1130; altarstone)|
The Corbridge garrison was composed of legionary cohorts taken at various times from several of the Roman legions which were stationed in Britain. The first to arrive was Legio II Augusta from Caerleon in south Wales, who were stationed here under governor Quintus Lollius Urbicus around AD140, followed in the late second century by cohorts of Legio VI Victrix from York. The cohorts from the Sixth Legion were augmented for a short time by contingents from Legio XX Valeria who were moved up from their legionary fortress at Chester in Cheshire during the administration of Sextus Calpurnius Agricola. The last legionary cohort recorded at Corbridge was from the Sixth, dated to the turn of the third century.
|CONCORDIAE LEG VI VI P F ET LEG XX|
|"To Concordia,¹ the Sixth Legion, Victorious, Loyal and Faithful and the Twentieth Legion [dedicates this]."|
Vexillus of the Second Legion Augusta
|LEG II AVG COH III F||"The Third Cohort of the Second Augustan Legion made this"||1155|
|LEG II AVG COH III||"The Third Cohort of the Second Augustan Legion"||1156|
|LEG II AVG COH IIII F||"The Fourth Cohort of the Second Augustan Legion made this"||1157|
|...CE LEG II AVG FEC||"[...] of the Second Augustan Legion made this"||1158|
The Second Legion is mentioned on at least eleven inscribed stones recovered from the Corbridge environs. Aside from the usual clutch of 'cohort stones' (vide supra) which proclaim responsibility for the structure of the defences and internal buildings, there are a couple of inscriptions which provide invaluable dating information (RIB 1147 & 1148), culturally important altarstones dedicated to classical gods (RIB 1127 & 1136), also a single tombstone to an unnamed soldier (RIB 1177); all texts shown below.
|DISCIPVLINAE AVGVSTORVM LEG II AVG|
|"For the Discipline of the Emperors, the Second Augustan Legion [made this]"|
|(RIB 1127; altar or statue base)|
|DEO SAN SILVANO MILITES VEXILLAT LEG II AVG ET COL LEGIVM SILVANIANO RVM ARAM DE SVO POS VOL LIB|
|"For the holy god Silvanus,¹ the soldiers from the Vexillation of the Second Augustan Legion and the College of the Silvaniani,² willingly and freely placed this altar out of their own resources."|
|(RIB 1136; altarstone)|
|IMP T AELIO ANIONINO AVG PIO II COS SVB CVRA Q LOLII VRBICI LEG AVG PR PR LEG II AVG F||IMP CAES I AELIO ANTONINO AVG PIO III COS P P SVB CVRA Q LOLLI VRBICI LEG AVG PR PR LEG II AVG FECIT|
|"For the emperor Titus Aelius Antoninus Augustus Pius, consul for the second time,¹ under the command of Quintus Lollius Urbicus,² legate of the Augustus with pro-praetorian power, the Second Legion Augusta built this"||"For Imperator Caesar Titus Aelius Antoninus Augustus Pius, three times consul,¹ Father of the Fatherland, under the care of Quintus Lollius Urbicus² the pro-praetorian legate of the emperor, the Second Augustan Legion made this"|
|(RIB 1147; dated: AD139)||(RIB 1148; dated: AD140)|
|D M MILES LEG II AVG ...|
|"To the spirit of the departed, a soldier of the Second Augustan Legion [...]."|
|(RIB 1177; tombstone)|
|LEG VI VIC FE||"The Sixth Victorious Legion made this"||1159>|
|LEG VI VIC P F||"The Sixth Victorious Legion, Loyal and Faithful"||1160|
|INSTANTE FL HYGIN > LEG VI V||"Restored by the century of Flavius Hyginus of the Sixth Victorious Legion"||1161|
|LEGIO VI PIE F VEX REFE||"The flag-section of the Sixth Legion, Loyal and Faithful rebuilt this"||1162|
There are at least a dozen stones bearing the name of the Sixth Legion; including three dateable to the latter half of the second century, three altarstones and a tombstone. This total includes the much-defaced and difficult inscription RIB 1190, which reads: ...IE... ...TITICIA... ...VI BRIV... ...TAE... ...L VI VIC ...F, the latter part of which contains the name of the legion.
|... VEXILLATIO LEG VI VIC P F SVB CN IVL VERO LEG AVG PER L O... TRIB MILITVM|
|"[...] flag-section of the Sixth Victorious Legion, Loyal and Faithful under Gnaeus Julius Verus,¹ legate of the emperor, through the agency of Lucius O[...] military tribune"|
|(RIB 1132; dated: c.AD158)|
|SOLI INVICTO VEXILLATIO LEG VI VIC P F F SVB CVRA SEX CALPVRNI AGRICO LAE LEG AVG PR PR|
|"To the Invincible Sun, the Vexillation of the Sixth Legion, Victorious, Loyal and Faithful, under the administration of Sextus Calpurnius Agricola,² pro-praetorian legate of the emperor"|
|(RIB 1137; dated: AD162-8)|
|VEXILLATIO LEG VI VIC P F SVB CVRA VIRI LVPI V C COS|
|"The Vexillation of the Sixth Legion, Victorious, Loyal and Faithful under the administration of the consular Virius Lupus,³ a most worthy man"|
|(RIB 1163; dated: AD197-202)|
|APOLLINI MAPONO Q TERENTIVS Q F OVF FIRMVS SAEN PRAEF CASTR LEG VI V P F D D||DEO MAPONO APOLLINI P AE... LVS > LEG VI VIC VSLM|
|"To Apollo Maponus, Quintus Terentius Firmus, son of Quintus, of the Aufentine voting tribe from Saena,¹ Praefectus Castrorum² of Legio Sextae Victrix Pia Fidelis, donated out of devotion"||"To the god Maponus Apollo, Publius Ae[lius Lucul]lus, centurion of the Sixth Victorious Legion willingly and deservedly fulfills his vow"|
|(RIB 1120; altarstone)||(RIB 1122; altarstone)|
|L VAL IVSTO MIL LEG VI EGN DYONISIVS ET SVR IVSTVS HER F C|
|"For Lucius Valerius Justus, a soldier of the Sixth Legion, Egn[atius] Dyonisius and Sur[ius] Justus his heirs had this made"|
|(RIB 1175; tombstone)|
|LEG XX V V FECIT||"The Twentieth Legion Valiant and Victorious made this"||1164|
|LEG XX V V FEC||"The Twentieth Legion Valiant and Victorious made this"||1165|
|LEG XXX¹ V V COH VII||"The Seventh Cohort of the Twentieth¹ Legion Valiant and Victorious"||1166|
The name of the Twentieth Legion appears on five stones from Roman Corbridge, and it seems that a Vexillatio of at least one cohort undertook building work during the latter part of the second century. This posting placed contingents from two separate legions in Corstopitum at the same time, and appears to have been cause of some disharmony, as the only altarstone dedicated by the Twentieth is one to Concordia, the dedication of which is significantly shared with the Sixth Legion (vide RIB 1125 supra).
|IMPERATORIBVS CAESARIBVS M AVRELIO ANTONINO AVG TRIBVNICIAE POTESTATIS ... COS ... ET L AVRELIO VERO AVG ARMENIACO TRIBVNICIAE POTESTATIS I... COS ...II VEXILLATIO LEG XX V V FECIT SVB CVRA SEXTI CALPVRNI AGRICOLAE LEGATI AVGVSTORVM PR PR|
|"For their imperial Caesars, Marcus Aurelius Antoninus Augustus, holding tribunician power for the [seventeenth]¹ time, consul [three]¹ times, and Lucius Aurelius Verus Augustus Armeniacus, holding tribunician power for the [third]¹ time, consul [two]¹ times, a Vexillation of the Twentieth Legion Valiant and Victorious made this, under the administration of Sextus Calpurnius Agricola, legate of the emperors with pro-praetorian power"|
|(RIB 1149; dated: AD163)|
|DIS MANIBVS FLAVINVS EQ ALAE PETR SIGNIFER TVR CANDIDI AN XXV STIP VII H S|
|"To the spirits of the departed and to Flavinus, a trooper of the Petrian Wing, standard-bearer in the turma of Candidus, twenty five years old with seven years service, here [he] lies"|
|(RIB 1172; tombstone)|
This cavalry unit is attested at Corbridge on at least one tombstone, and is possibly mentioned on another, both stones being shown here. The wing was named after Titus Pomponius Petra, a former commander, and the name was later transferred to their main garrison fort Petrianvm (Stanwix, Cumbria). For further information see Ala Petriana.
|...MERITO EX EQ ALAE ...AE|
|"[...] veteran and former trooper of the Wing [...]"|
|(RIB 1178; suspected tombstone)|
|DISCIP AVGVSTO MILITES COH I F VARDVLLORVM M C R EQ CVI PRAEEST PVB CALPVRNIVS VICTOR TR|
|"For the Discipline of the Emperor, the soldiers of the First Cohort of Vardulli, one-thousand strong, citizens of Rome, part-mounted, under the command of the tribune Publius Calpurnius Victor [made this]."|
The First Cohort of Vardulli is attested at Corbridge on a single dedicatory stone (see text above). This part-mounted unit were originally recruited from among the Vardulli tribe, who inhabited Hispania Terraconensis, Guipuscoa, northern Spain. The unit is also recorded on two stones from Longovicium (Lanchester, Durham; RIB 1076 & 1083), on seven inscriptions variously dated between AD216 to AD241 from Bremenium (High Rochester, Northumberland; RIB 1262, 1263, 1272, 1279, 1280, 1285 & 1288), and on single undated altarstones in Scotland at Cappuck in Borders Region (RIB 2118) and Castlecary in Central (RIB 2149).
|COH [I] LING ILIOMARVS|
|"The [First] Cohort of Lingones, [century of] Juliomarus [made this]"|
Cohors Primae Lingonum Equitata was a part-mounted unit originating from the Lingones tribe of Gallia Lugdunensis, inhabiting the Bourgogne region of Central France. The First Cohort of Lingones is known from inscriptions at Bremenium (High Rochester, Northumberland; RIB 1276; AD139-43) and Longovicium (Lanchester, Durham; RIB 1091/1092; AD238-44), and possibly also here at Corbridge, where it is recorded on a single undated stone (RIB 1186 supra), which is missing the unit number.
This large early installation is described on a separate page: Beaufront Red House.
A Trajanic coin sealed beneath the rampart of the Stanegate fort at Corbridge proves a foundation date of AD103 or later, at the same time that the emperor Trajan was withdrawing troops from Britain (and elsewhere) for deployment in his second Dacian campaign which commenced in 105. It would appear that the original (Agricolan?) fort was burnt to the ground and the area levelled shortly after AD103. This is not indicative of barbarian activity, who would hardly be mindful to carefully level the ground after a night's arson attack, but is sure evidence of careful preparation of the foundations for another, larger fort built upon the same site during the early-Trajanic period. This complete rebuilding of the former fort is evidence of a change in the type of garrison unit housed at the site.
There is evidence of another rebuild and accompanying change of garrison in the early Hadrianic period, and a certain amount of rebuilding during late-Antonine times in the mid-second century is attested by an inscription of governor Sextus Calpurnius Agricola, who replaced Priscus around AD162. This probably indicates the strengthening of the Hadrianic and Stanegate barriers following the withdrawal from the Antonine Wall in Scotland. By the 3rd century Corbridge had grown into a large sprawling garrison town of 12 hectares enclosed by walls and housing a legionary garrison at its centre.
The dimensions of the auxiliary forts at Corbridge are unknown, but judging from the size of the garrison units, two of which contained a nominal one-thousand men and/or horses, the forts must have covered an area of at least 8 or 9 acres (c.3.6ha). It is noteable that of all the forts on the Northern frontier only Corbridge has yielded Saxon artefacts, perhaps indicating that the Hadrianic barrier a couple of miles to the north continued to keep barbary at bay, at least for a while.
|IMP CAES L SEP SEVERVS PI PERTINAX ET IMP CAESAR M AVR ANTONINVS PIVS AVGVSTI ET P SEPTIMIVS [GETA] CAESAR HORREVM PER VEXILLATIONEM LEG ... FECERVNT SVB L ALFENO SENECIONE LEG AVGG PR PR|
|"For Imperator Caesar Lucius Septimius Severus Pius Pertinax, Imperator Caesar Marcus Aurelius Antoninus Pius Augustus and Publius Septimius Geta Caesar, this granary, through the agency of a detachment of the [...] Legion, was built under Lucius Alfenus Senecio,¹ legate of the emperors with pro-praetorian power."|
|(RIB 1151; restored inscription; dated: AD205-8)|
The latest datable inscription found at Corbridge is RIB 1151, a restored inscription for which there are two feasible concluding lines (see above). Lucius Alfenus Senecio governed Britain from c.AD205 to 208 and is the last known governor of the entire British province before it was partitioned by the emperor Septimius Severus sometime before his death at York in 211. Gaius Valerius Pudens was the immediate predecessor of Senecio, and governed Britain from c.AD202 to 205.
During excavations over the years at Corbridge a number of animal bones have been uncovered, including those of domesticated Ox, Sheep, Goats and Pigs, game such as Red Deer, Roe Deer, Wild Ox and Hare, also animals such as Fox, Badger, Beaver, Vole and Mole; the latter group very likely being hunted and killed for sport and as a means of pest control. Among the bones recovered from the Red House site were those of Ox, Sheep, Goat, Pig, Red Deer and Roe Deer.
"At Corbridge, records have been obtained of roads and buildings over a wide area around the visible remains exposed by excavation. The main street fronting the two military compounds continues in an irregular course east and west. It is flanked by buildings, and other streets branch off to north and south. A third of a mile to the west is a Romano-Celtic temple of normal plan. The outline of the precinct-wall, which encloses an area of perhaps 120 by 110 ft., and of a central building are clearly visible, though no trace remains on the surface." (St. Joseph, 1951)
|D M ...RATHES PAL MORENVS VEXILA VIXIT AN LXVIII|
|"To the spirits of the departed [and Aria]rathes¹ Morenus the Palmyrene,² vexillarius² who lived for sixty-eight years."|
|(RIB 1171; tombstone)|
As always, the best epigraphic evidences of civilian settlement at Corbridge come in the form of tombstones.
|D M IVL PRIMVS CONIVGI C P C||"To the spirits of the departed and to Julius Primus, husband. His wife placed this as arranged."||1174|
|D M AHTEHE FIL NOBILIS VIXIT ANIS V||"To the spirits of the departed and Athene, an excellent daughter who lived for five years."||1180|
|D M SVDRENVS ERTOLE NOMINE VELLIBIA FELICISSIME VIXIT ANIS IIII DIEBVS LX||"To the spirits of the departed, Sudrenus Ertole nominates the most happy Vellibia, who lived for four years and sixty days."||1181|
|IVLIA MATER NA AN VI IVL MARCELLINVS FILIAE CARISSIMAE||"Julia Materna, six years old. Julius Marcellinus [made this] for a most lovely daughter."||1182|
|LEG A... Q CALPVRNIVS CONCESSINI VS PRAEF EQ CAESA CORIONOTOTARVM MANV PRAESENTISSIMI NVMINIS DEI V S|
|"The Legate of the Augustus [...] for cutting-down an armed band of Corionototae,¹ Quintus Calpurnius Concessinius, Prefect of Cavalry, fulfills his vow to the spirit of the most omnipresent god.²."|
|(RIB 1142; altarstone)|
Over twenty altarstones have been uncovered at Corbridge, mostly dedicated to various gods from the classical pantheons of Greece and Rome, although the greatest number of altars to a single god is the Romano-Celtic amalgam Apollo Maponus, to whom there are four dedications, closely followed by the Germanic god Veterus with three. The only other gods possessing more than one dedication are Jupiter and Discipline, each with two altarstones. There are single altars dedicated to Astarte (in Greek), Concordia, Diana, Hercules (in Greek), Mercury (in relief), Minerva, Panthea, Silvanus, Sol Invictus (Mithras) and Victory; there are another six altarstones dedicated to gods whose names are illegible or otherwise unknown. A selection of the more interesting examples are shown here.
|IOVI AETERNO DOLICHENO ET CAELESTI BRIGANTIAE ET SALVTI G IVLIVS APOLINARIS > LEG VI IVSS DEI|
|"To the eternal Jupiter of Doliche,¹ Celestial Brigantia² and Salus,³ Gaius Julius Apolinaris, Centurion of the Sixth Legion, [set this up] by command of the god."|
|(RIB 1131; altarstone)|
|ASTARTES BOMON M HESORAS POULCHER M ANETHEKEN||HERAKLEI TYRIOI DIODORA ARCHIERIA|
|"You see me, an altar of Astarte,¹ Pulcher set me up."||"To Herakles of Tyre,² the priestess Diodora (set this up)."|
|(RIB 1124; altarstone; in Greek)||(RIB 1129; altarstone; in Greek)|
|VICTORIAE AVG L IVL IVLIANVS LEG AVG ...VS...|
|"To August Victory, Lucius Julius Julianus, legate of the emperor [... willingly and deservedly] fulfilled his vow."|
|(RIB 1138; altarstone)|
|DEO VETERI||VIT M ITI||DEO VITIRI|
|"For the god Veterus."||"For Vitiris, Marcus Itius [made this]."||"To the god Vitiris."|
|(RIB 1139; altarstone)||(RIB 1141; altarstone)||(RIB 1140; altarstone)|
There are three altarstones recovered from Corbridge which are dedicated to the god Veterus or Vitiris (vide supra), an ancient German ancestral deity worshipped in Britain under a variety of names including; Veter, Veteres, Viter and Votris. The god is also known from altars at Concangis (Chester-le-Street, Durham; RIB 1046), Vindomora (Ebchester, Durham; RIB 1103) and Cataractonium (Catterick, North Yorkshire; RIB 727), also at many forts along Hadrian's Wall.
|DEAE MINERVAE T TERTINIVS... LIBR EX VOTO POS|
|"To the goddess Minerva,¹ Titus Tertinius [...] Librarius,² placed this as the result of a vow."|
|(RIB 1134; altar or statue base)|
|APOLLINI MAPONO CALPVRNIVS ... TRIB DEDICAVIT||DEO ARECVRIO APOLLINARIS CASSI VSLM|
|"To Apollo Maponus, the tribune Calpurnius [...] has dedicated [this]."||"To the god Arecurius Apollo, Cassius willingly and deservedly fulfills his vow."|
|(RIB 1121; altarstone)||(RIB 1123; altarstone)|
Apollo, also known by the Romans as Phoebus (the sun), was the son of Jupiter and Latona, and brother of Diana (a.k.a. Phoebe, the moon). He was the god of the fine arts, music, poetry, medicine and eloquence, and reputed to be master of the bow and arrow, as was his sibling goddess. His temples are known throughout the Roman world, including many examples in Britain.
|ARA DIAN POSVI N...||DEO MERCVRIO||B F DEAE PANTHEAE|
|"An Altar for Diana,¹ placed from us [...]."||"To the god Mercury.²"||"Good fortune to the goddess Panthea.³"|
|(RIB 1126; altarstone)||(RIB 1133; relief of Mercury)||(RIB 1135; altarstone)|
|... ... ...SIT... ... ...NORVS ...PRAEP CVRAM AGENS HORR TEMPO RE EXPEDITIONIS FELICISSI BRITTANNIC VSLM|
|"[...] the acting administrator,¹ planning in advance [so that] the granaries were repaired in time for the successful campaigns in Britain, willingly and deservedly fulfilling a vow."|
|(RIB 1143; altarstone)|
|IMP C M PIAVONIO VICTORINO P F AVG||AVG ... CAESAR MAXIMINVS AVG N|
|"Imperator Caesar Marcus Piavonius Victorinus Pius Felix Augustus.¹"||"Augustus [...]² Caesar Maximinus, our emperor.³"|
|(RIB 2296; dated: AD269-271)||(RIB 2297; dated: AD235-238)|
The sanctuary area of Corstopitum lay in two sections to the north of the military enclosures at the heart of the Roman town. The defenses of both the eastern and western compounds have a very un-military outline due to their methodical respect of the temples sacred boundaries. All of the temples so far discovered appear to be constructed in the classical style, which is to be expected in a town with a predominantly legionary population, all of whom were Roman citizens and thus inclined towards the classical pantheon. The eastern enclave contains at least five known temples (numbers 1 through 5) while the western enclave holds two (6 & 7). Unfortunately, although there have been several altarstones and religious artifacts turned-up in Corbridge over the years, none may be positively assigned to any of the classical shrines.
The podium of this temple was composed of packed earth held within retaining walls of dressed stone measuring 24½ ft. wide by 33 ft. long. There were five irregularly-spaced columns along the northern front, the bases of which were 1 ft 4 ins. square, which would imply a column-height of between 10-12 ft. The north-east corner of the temple was destroyed, possibly during barbarian incursions south of Hadrian's Wall around AD296.
Built alongside Temple 1 only 2 ft. to the east, the podium of this temple measured 31 ft. 5 ins. wide and at least 55½ ft. long. There were four columns set along the front of the temple, two spaced 10 ft. apart to either side of the door leading to the cella, the sanctuary of the temple. The form of the temple sanctuary was an open courtyard with surrounding roofed colonnade containing a massive platform set at the rear, probably to house an altar which was also open to the sky.
Lying to the immediate east of Temple 2, all that survives of this temple is the front of the podium measuring 27 ft. wide; the original length is unknown.
This temple is situated to the north-east of Temple 3, behind Temple 5, and unlike the three preceeding temples faced either west or east. Only the podium has survived, measuring 27 ft. 3 ins. wide by 32 ft. 8 ins. long.
Like temple 4, this temple is oriented east-west and is known only from its podium, which measured about 26½ ft. wide by at least 43 ft. in length. It was situated to the immediate north of Temple 3 and just west of Temple 4, obscuring them both.
This temple lies in the western enclave and is known only from its podium, which measures 12 ft. 8 ins. in width by 24 ft. 10 ins. long. It is the smallest known temple at Corbridge and is oriented north-south, its facade probably opened onto the street to the north.
This temple lies to the immediate south of Temple 6 and is known only from the south-east corner of its podium, which was probably aligned east-west. |
Altitude or height is defined based on the context in which it is used (aviation, geometry, geographical survey, sport, and more). As a general definition, altitude is a distance measurement, usually in the vertical or "up" direction, between a reference datum and a point or object. The reference datum also often varies according to the context. Although the term altitude is commonly used to mean the height above sea level of a location, in geography the term elevation is often preferred for this usage.
Vertical distance measurements in the "down" direction are commonly referred to as depth.
Altitude in aviation and in spaceflight
In aviation, the term altitude can have several meanings, and is always qualified by either explicitly adding a modifier (e.g. "true altitude"), or implicitly through the context of the communication. Parties exchanging altitude information must be clear which definition is being used.
Aviation altitude is measured using either Mean Sea Level (MSL) or local ground level (Above Ground Level, or AGL) as the reference datum.
Pressure altitude divided by 100 feet (30m) as the flight level, and is used above the transition altitude (18,000 feet (5,500 m) in the US, but may be as low as 3,000 feet (910 m) in other jurisdictions); so when the altimeter reads 18,000 ft on the standard pressure setting the aircraft is said to be at "Flight level 180". When flying at a Flight Level, the altimeter is always set to standard pressure (29.92 inHg / 1013.25 mbar).
On the flight deck, the definitive instrument for measuring altitude is the pressure altimeter, which is an aneroid barometer with a front face indicating distance (feet or metres) instead of atmospheric pressure.
There are several types of aviation altitude:
- Indicated altitude is the reading on the altimeter when the altimeter is set to the local barometric pressure at Mean Sea Level.
- Absolute altitude is the height of the aircraft above the terrain over which it is flying. Also referred to feet/metres above ground level (AGL).
- True altitude is the actual elevation above mean sea level. It is Indicated Altitude corrected for non-standard temperature and pressure. In UK aviation radiotelephony usage, the vertical distance of a level, a point or an object considered as a point, measured from mean sea level; this is referred to over the radio as altitude.(see QNH)
- Height is the elevation above a ground reference point, commonly the terrain elevation. In UK aviation radiotelephony usage, the vertical distance of a level, a point or an object considered as a point, measured from a specified datum; this is referred to over the radio as height, where the specified datum is the airfield elevation (see QFE)
- Pressure altitude is the elevation above a standard datum air-pressure plane (typically, 1013.25 millibars or 29.92" Hg). Pressure altitude and indicated altitude are the same when the altimeter is set to 29.92" Hg or 1013.25 millibars.
- Density altitude is the altitude corrected for non-ISA International Standard Atmosphere atmospheric conditions. Aircraft performance depends on density altitude, which is affected by barometric pressure, humidity and temperature. On a very hot day, density altitude at an airport (especially one at a high elevation) may be so high as to preclude takeoff, particularly for helicopters or a heavily loaded aircraft.
These types of altitude can be explained more simply as various ways of measuring the altitude:
- Indicated altitude – the altimeter reading
- Absolute altitude – altitude in terms of the distance above the ground directly below
- True altitude – altitude in terms of elevation above sea level
- Height – altitude in terms of the distance above a certain point
- Pressure altitude – the air pressure in terms of altitude in the International Standard Atmosphere
- Density altitude – the density of the air in terms of altitude in the International Standard Atmosphere
- Troposphere — surface to 8,000 metres (5.0 mi) at the poles – 18,000 metres (11 mi) at the equator, ending at the Tropopause.
- Stratosphere — Troposphere to 50 kilometres (31 mi)
- Mesosphere — Stratosphere to 85 kilometres (53 mi)
- Thermosphere — Mesosphere to 675 kilometres (419 mi)
- Exosphere — Thermosphere to 10,000 kilometres (6,200 mi)
High altitude and low air pressure
Regions on the Earth's surface (or in its atmosphere) that are high above mean sea level are referred to as high altitude. High altitude is sometimes defined to begin at 2,400 metres (8,000 ft) above sea level.
At high altitude, atmospheric pressure is lower than that at sea level. This is due to two competing physical effects: gravity, which causes the air to be as close as possible to the ground; and the heat content of the air, which causes the molecules to bounce off each other and expand.
Because of the lower pressure, the air expands as it rises, which causes it to cool. Thus, high altitude air is cold, which causes a characteristic alpine climate. This climate dramatically affects the ecology at high altitude.
Relation between temperature and altitude in Earth's atmosphere
The environmental lapse rate (ELR), is the rate of decrease of temperature with altitude in the stationary atmosphere at a given time and location. As an average, the International Civil Aviation Organization (ICAO) defines an international standard atmosphere (ISA) with a temperature lapse rate of 6.49 K(°C)/1,000 m (3.56 °F or 1.98 K(°C)/1,000 Ft) from sea level to 11 kilometres (36,000 ft). From 11 to 20 kilometres (36,000 to 66,000 ft), the constant temperature is −56.5 °C (−69.7 °F), which is the lowest assumed temperature in the ISA. The standard atmosphere contains no moisture. Unlike the idealized ISA, the temperature of the actual atmosphere does not always fall at a uniform rate with height. For example, there can be an inversion layer in which the temperature increases with height.
Effects of high altitude on humans
Medicine recognizes that altitudes above 1,500 metres (4,900 ft) start to affect humans, and extreme altitudes above 5,500–6,000 metres (18,000–20,000 ft) cannot be permanently tolerated by humans. As the altitude increases, atmospheric pressure decreases, which affects humans by reducing the partial pressure of oxygen. The lack of oxygen above 2,400 metres (8,000 ft) can cause serious illnesses such as altitude sickness, high altitude pulmonary edema, and high altitude cerebral edema. The higher the altitude, the more likely are serious effects. The human body can adapt to high altitude by breathing faster, having a higher heart rate, and adjusting its blood chemistry. It can take days or weeks to adapt to high altitude. However, above 8,000 metres (26,000 ft), (in the "death zone"), the human body cannot adapt and will eventually die.
There is a significantly lower overall mortality rate for permanent residents at higher altitudes. Additionally, there is a dose response relationship between increasing elevation and decreasing obesity prevalence in the United States. However, people living at higher elevations have a statistically significant higher rate of suicide. The cause for the increased suicide risk is unknown so far.
For athletes, high altitude produces two contradictory effects on performance. For explosive events (sprints up to 400 metres, long jump, triple jump) the reduction in atmospheric pressure signifies less atmospheric resistance, which generally results in improved athletic performance. For endurance events (races of 5,000 metres or more) the predominant effect is the reduction in oxygen which generally reduces the athlete's performance at high altitude. Sports organisations acknowledge the effects of altitude on performance: the International Association of Athletic Federations (IAAF), for example, have ruled that performances achieved at an altitude greater than 1,000 metres (3,300 ft) will not be approved for record purposes.
Athletes also can take advantage of altitude acclimatization to increase their performance. The same changes that help the body cope with high altitude increase performance back at sea level. These changes are the basis of altitude training which forms an integral part of the training of athletes in a number of endurance sports including track and field, distance running, triathlon, cycling and swimming.
Effect of altitude on animals
Decreased oxygen availability and decreased temperature make life at high altitude challenging. Despite these environmental conditions, many species have been successfully adapted at high altitudes. Animals have developed physiological adaptations to enhance oxygen uptake and delivery to tissues which can be used to sustain metabolism. The strategies used by animals to adapt to high altitude depend on their morphology and phylogeny.
Fish at high altitudes may also have a lower metabolic rate, as has been shown in highland westslope cutthroat trout compared to introduced lowland rainbow trout in the Oldman River basin. There is also a general trend of smaller body sizes and lower species richness at high altitudes observed in aquatic invertebrates, likely due to lower oxygen partial pressures. These factors may decrease productivity in high altitude habitats, meaning there will be less energy available for consumption, growth, and activity, which provides an advantage to fish with lower metabolic demands.
The naked carp from Lake Qinghai, like other members of the carp family, can use gill remodelling to increase oxygen uptake in hypoxia. The response of naked carp to cold and low-oxygen conditions seem to be at least partly mediated by hypoxia-inducible factor 1 (HIF-1). It is unclear whether this is a common characteristic in other high altitude dwelling fish or if gill remodelling and HIF-1 use for cold adaptation are limited to carp.
Rodents living at high altitude include deer mice, guinea pigs and rats. As small mammals they face the challenge of maintaining body heat in cold temperatures, due to their large volume to surface area ratio. As oxygen is used as a source of metabolic heat production, the hypobaric hypoxia at high altitudes is problematic.
There are a number of mechanisms that help them survive these harsh conditions including altered genetics of the hemoglobin gene in guinea pigs and deer mice. Deer mice use a high percentage of fats as metabolic fuel at high altitude to retain carbohydrates for small burst of energy. To convert fats to energy in the form of ATP, more oxygen is required than to convert the same amount of carbohydrates. The reason they use fats is believed to be because they have it in large stores, but also means that they must eat more or they will begin to lose weight.
Other physiological changes that occur in rodents at high altitude include increased breathing rate and altered morphology of the lungs and heart allowing more efficient gas exchange and delivery. Lungs of high altitude mice are larger, with more capillaries, and hearts of mice and rats at high altitude have a heavier right ventricle, which pumps blood to the lungs.
Birds have been especially successful at living at high altitudes. In general, birds have physiological features that are advantageous for high-altitude flight. The respiratory system of birds moves oxygen across the pulmonary surface during both inhalation and exhalation, making it more efficient than that of mammals. In addition, the air circulates in one direction through the parabronchioles in the lungs. Parabronchioles are oriented perpendicular to the pulmonary arteries, forming a cross-current gas exchanger. This arrangement allows for more oxygen to be extracted compared to mammalian concurrent gas exchange; as oxygen diffuses down its concentration gradient and the air gradually becomes more deoxygenated, the pulmonary arteries are still able to extract oxygen. Birds also have a high capacity for oxygen delivery to the tissues because they have larger hearts and cardiac stroke volume (mL / min) compared to mammals of similar body size. Additionally, they have an increased vascularization in flight muscle due to increased branching of capillaries and small muscle fibres (which increases surface-area-to-volume ratio). These two features facilitate oxygen diffusion from the blood to muscle, allowing flight to be sustained during environmental hypoxia. Bird's hearts and brains, which are very sensitive to arterial hypoxia, are more vascularized compared to mammals. The bar-headed goose (Anser indicus) is an iconic high flyer that surmounts the Himalayas during migration, and serves as a model system for derived physiological adaptations for high-altitude flight.
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- Downloadable ETOPO2 Raw Data Database (2 minute grid)
- Downloadable ETOPO5 Raw Data Database (5 minute grid)
- Find the altitude of any place
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Many programs are designed to enhance physicians' communication skills, but less attention has been devoted to helping patients improve their communication during medical visits. Nevertheless, this small body of work has grown substantially since the earliest patient activation interventions were begun some 30 years ago (1,2).
In those studies, an intervention conducted in the waiting room to increase patient participation in the medical dialogue was successful in changing patients' communication behavior. A systematic review of patient activation interventions through 2004 identified some 20 studies of this kind, the majority of which reported modest but significant effects, most often in question asking (3).
None of the activation interventions in the 2004 review took place within the context of mental health care, although it is well documented that psychiatric patients often have expectations and preferences for care but do not readily communicate them to their psychiatrists, who in turn fail to identify what their patients want (4,5). In a recent study, Alegráa and colleagues (6) developed and evaluated an activation and empowerment intervention designed to enhance question-asking skills and self-management strategies for patients receiving care in community mental health clinics. The three-session intervention was successful in increasing patients' self-reported activation during subsequent mental health visits, including greater engagement in information exchange and retention in treatment.
The promising results of that study suggest that psychiatric patients, like others, can benefit from an activation intervention by changing communication during visits to therapists and enhancing their commitment to continued care. However, it appears that no other studies have focused on educating mental health patients about evidence-based treatment, providing feedback, and encouraging discussion with their clinician.
This study furthered this line of work by developing an Internet-based interactive computer tool to educate patients with schizophrenia about evidence-based treatment guidelines and to compare their care to guidelines and receive feedback (7). The design of the intervention was also motivated by a study that found disparities between the evidence-based treatment recommendations of the Schizophrenia Patient Outcomes Research Team, known as PORT, and the services received by persons with schizophrenia (8,9). The study found substantial underuse of psychosocial treatments, including employment and family education interventions, undertreatment of adjunctive depression, and underuse of prophylactic antiparkinson medication. Improving patients' abilities to participate actively in their own treatment and management may lead to increased use of such services.
The goal of the Web-based intervention was to empower patients with schizophrenia to discuss their mental health treatment with their therapist. Patients were randomly assigned to the intervention or to a control group. Subsequently visits to clinicians by each group of patients were audiorecorded and compared by using the Roter Interaction Analysis System (RIAS), which characterizes medical dialogue. The hypotheses were that patients assigned to the intervention would be more verbally active in their visits, particularly in regard to asking questions about drug therapy and side effects and discussing psychosocial and lifestyle issues, and that clinicians' communications with patients who completed the intervention would be more patient-centered than their communications with the other patients.
All patients and clinicians gave signed informed consent to participate. The study was reviewed and approved by the Johns Hopkins Bloomberg School of Public Health Institutional Review Board.
Study sites and clinician participants
The three study sites were an assertive community treatment (ACT) program and two hospital psychiatric outpatient clinics. A total of 68 clinicians, including psychiatrists, social workers, psychologists, nurses, and counselors, were identified as eligible for the study. Because psychiatrists at the ACT site primarily saw patients for brief medication management visits, only the six nonpsychiatrists at the site were considered eligible. All clinicians at the two psychiatric outpatient clinics, 20 psychiatrists and 42 clinicians from other disciplines, were eligible for recruitment.
Twenty clinicians participated in the study by referring one or more patients with schizophrenia. Participants included five of the 20 (25%) psychiatrists and nine of the 42 (21%) other clinicians at the outpatient clinics and all six nonpsychiatrist clinicians at the ACT site.
Patient eligibility included referral by a study clinician and a scheduled visit during the 18-month study period. Additional criteria were age of 13 or older, a diagnosis of schizophrenia, literacy in English, and ability to give informed consent. A total of 163 eligible patients were identified by research staff through patient lists or records. The refusal rate, either because a clinician declined to refer the patient or the patient directly refused to participate, averaged 55% at the two outpatient clinics and was 30% at the ACT program, resulting in enrollment of 56 patients. Patients were assigned at random to the intervention or a control group.
YourSchizophreniaCare is an interactive Web site designed by the research team and programmed first by the Foundation for Accountability and later by HealthGrades, Inc., with CompareYourCare software. YourSchizophreniaCare is a patient-oriented, Web-based learning approach to help patients navigate six areas of quality of care—medications, side effects, referrals, family support, employment, and quality of life. Its goal was to increase the likelihood that patients will initiate discussion of these topics with their therapist. The average time to complete the intervention was 20 minutes.
For each area, patients answer questions about their current status and treatment. On the basis of the pattern of patient responses, individualized feedback recommendations appear on the screen. For instance, patients are asked how often medicine makes them feel restless or jittery inside. If they report medication side effects, they are encouraged to discuss them with their therapist at an upcoming visit.
The feedback is supplemented with video clips of an actor simulating a patient as he or she implements a recommendation for discussing a specific topic with the doctor. The Web site includes 14 unique 30-second video clips designed to model the performance of communication strategies and skills. The clips provide examples of how someone with schizophrenia can broach potentially sensitive issues with a therapist, such as confusion about prescribed drugs, side effects and poor adherence, use of alcohol with medication, barriers to more consistent treatment adherence, and family tensions. Communication strategies to assist the patient to be proactive were also modeled in video clips. They included setting the visit agenda by expressing expectations and goals, using paraphrase and interpretation to check for understanding, asking for plain words instead of scientific terms, and summarizing important information at the close of the visit.
The control group viewed a 22-minute video about schizophrenia treatment prepared by the Mental Illness Research, Education, and Clinical Center at the U.S. Department of Veterans Affairs Medical Center in Baltimore (10). They also received a brochure written in lay language that presents five schizophrenia treatment recommendations addressed by the Web-based tool, including maintenance dose of antipsychotic medication, prophylactic antiparkinson medication, antidepressant treatment, family intervention, and supported employment (7). The combination of brochure and video mirrored the information content of the Web-based intervention but without the personalized feedback or video clips role-modeling discussion of treatment issues provided online.
The generalized estimating equation method for correlated data was used in all regression analyses in order to account for nonindependence across observations (11). An exchangeable correlation structure was assumed by using robust estimation, which is likely to yield more accurate or valid coefficient estimates, even if the correlation structure is misspecified (12). Because some clinicians practiced at more than one site, analyses were not designed to account for intraclass correlation within sites. Covariates included clinician training status (psychiatrist or nonpsychiatrist), clinician gender, patient education level, and patient gender. Other covariates, including race-ethnicity, were not statistically significant in bivariate analyses and were not included in the multivariate model.
Audiorecordings of patient visits were coded with RIAS, a widely used system for characterizing medical dialogue with well-established reliability and predictive validity (13). Each statement by patient and clinician, defined as a complete thought, is assigned to mutually exclusive and exhaustive codes. Statements are coded directly from recordings without transcription. Several RIAS categories address biomedically focused behaviors, including information giving, question asking, and counseling about general medical and psychiatric symptoms and history and current treatment and regimens, including medication and appointments. Lifestyle and psychosocial categories capture exchange related to the patient's living situation, familial relations, and functions related to daily living or emotional states such as depression or anxiety. Psychosocial exchange also includes any discussion of alcoholism or drug abuse. Socioemotional communication categories for positive or negative exchange and responses to emotion such as empathy, concern, and legitimation are also coded.
Structural indicators of communication include visit duration in minutes, the sum of all patient and clinician statements, and a measure of clinician verbal dominance constructed as the ratio of all clinician statements—excluding such “back channels” as “Go on,” “Yes,” Uh huh, and “Right”—to patient statements. A value greater than 1.0 indicates that the clinician contributed more than the patient to medical dialogue.
In this study, as in other studies that used the RIAS system, patient-centered communication was operationalized as proportionately greater engagement by the patient in information seeking and disclosure in the psychosocial and socioemotional dimensions of illness management than in biomedically focused dialogue about disease management (14,15). The patient-centeredness ratio was calculated as the number of patient-centered statements divided by the number of medically focused and directive statements. Specifically, patient-centered statements included patients' psychosocial and medical questions, emotional statements, and information plus the clinicians' psychosocial-related questions, statements providing psychosocial and medical information, statements related to emotional topics, and partnership building statements. Medically focused and directive statements included clinicians' medical questions and statements about procedures and directions and patients' statements providing biomedical information.
Coders also globally rated the emotional tone of the patient and the clinician across several dimensions, such as warmth and sympathy, engagement, friendliness, dominance, hurry, and anxiety, on a scale of 1, low or none, to 5, high. The two coders were not aware of study hypotheses or patients' intervention status. A random sample of five audiotapes was drawn throughout the coding period for double coding to establish intercoder reliability. Pearson correlation coefficients between random pairs of coders were calculated for all categories with an average frequency greater than one. Reliability averaged .88 for provider categories (range .61–.99) and .89 for patient categories (range .81–.99) and was similar to reliability reported by other RIAS studies (13). For global affect ratings, coders were in 100% agreement within one scale point for provider ratings and averaged 98% agreement within one scale point for patient ratings.
Technical failures in audiorecordings reduced the sample sizes from 56 to 50 patients and from 20 to 19 clinicians. Clinicians included five psychiatrists, five licensed clinical professional counselors, one licensed graduate professional counselor, four licensed graduate social workers, two licensed clinical social workers, a clinical psychologist, a registered nurse with a master's degree in nursing, and another clinician with unspecified training. All the counselors and social workers had master's degrees at minimum; differences in titles reflect Maryland state licensure requirements.
The average number of patients enrolled in the study for each clinician was 2.8 (range one to 12). There were no statistically significant differences in the number of patients in the intervention and control groups seen by psychiatrists or other clinicians. Eleven of the visits by patients in the intervention group were with psychiatrists, five were with clinical counselors, five were with clinical social workers, one was with a clinical psychologist, one was with a registered nurse with a master's degree in nursing, and one was with a clinician with unspecified training. Fifteen of the visits by patients in the control group were with psychiatrists, one was with a clinical counselor, seven were with clinical social workers, two were with clinical psychologists, and one was with a clinician with unspecified training.
Overall, patients were predominantly African American (N=34, 68%), had less than a high school education (N=29, 58%), had an average age of 50, and had never married (N=37, 74%). A total of 16 (32%) were living with family members, 16 (32%) were not living with family but saw them regularly, and 18 (36%) did not see family regularly or had no family in the area. Most patients (N=44, 88%) obtained health insurance coverage through Medicare, Medicaid, or the Department of Veterans Affairs. Fifteen (30%) patients reported receiving Social Security Disability Insurance, and 28 (56%) received Supplemental Security Income. Eight (16%) patients reported some earned income, and six (12%) reported receiving regular contributions from family and other sources. There were no statistically significant differences between intervention and control groups on any of these measures (Table 1).
As shown in Table 2, visits by patients in the intervention and control groups differed in a number of ways. Visits by patients in the intervention group were several minutes longer compared with visits by patients in the control group (23.7 versus 19.3 minutes, respectively; p<.05). Patients in the intervention group contributed more actively to the dialogue compared with patients in the control group (288 versus 229 statements per visit; p<.05). The number of statements by clinicians per visit and the total number of statements per visit were higher among visits by the intervention group compared with visits by the control group, but the differences were not statistically significant. Consistent with these dialogue characteristics, the verbal dominance of the clinician was lower in the intervention visits compared with control group visits.
The patient-centeredness ratio was significantly higher during visits by patients in the intervention group compared with visits by patients in the control group (8.5 versus 3.2, p<.05), as shown in Table 2. A higher patient-centeredness ratio indicates that the visit featured a relatively greater emphasis on psychosocial and lifestyle categories in contrast to biomedical exchange.
More specifically, patients in the intervention group asked more questions about their treatment regimen and tended to discuss psychosocial issues more often than control group patients. They also gave more information about their lifestyles to the therapist and were more likely to check their understanding of information communicated to them (Table 3). There were no differences between the groups in the number of socioemotional statements, such as expressions of concern, optimism, agreement, or disagreements (data not shown).
During visits with patients in the intervention group, clinicians were significantly less likely to ask questions about the treatment regimen or medical symptoms but were more likely to express empathy and to use more facilitators and cues of interest to encourage the patient to continue to speak (Table 3). There were no differences between visits by the two groups of patients in other categories of clinician socioemotional or biomedical exchange (data not shown).
Emotional tone of the dialogue differed for both patients and clinicians depending on the study group (Table 4). Patients in the intervention group were rated as more dominant and respectful than those in the control group. They were also rated as sounding more distressed, but the trend did not reach statistical significance. During visits with patients in the intervention group, clinicians were rated as sounding more sympathetic and warm than during visits with the control group; they were also rated as sounding more engaged in dialogue with intervention group patients, but the trend was not statistically significant.
The patient activation intervention was successful in achieving its primary goals. As hypothesized, patients exposed to the intervention were more verbally active during mental health visits than control group patients. Visits with intervention group patients were longer by several minutes, and patients in that group contributed proportionately more to the medical dialogue than control patients. As a result, clinicians were less verbally dominant during visits with patients exposed to the Web site.
In addition, differences between groups in several specific elements of communication were evident. Patients enrolled in the intervention asked more questions about their therapeutic regimen and psychosocial and lifestyle issues and gave more information to the clinician about lifestyle issues. They also were more likely to check that they understood what the clinician said by restating or paraphrasing information to confirm a shared understanding of the facts or of the issues being discussed. In addition, patients in the intervention group were rated as sounding more dominant and respectful in their communication but also as more distressed.
In a somewhat parallel manner, clinician communication also showed intervention effects. During visits with patients in the intervention group, clinicians encouraged patient disclosure by higher use of facilitators and cues of interest and responded more empathically. They also asked fewer questions about the treatment regimen and mental disorder symptoms than during visits with control group patients. This pattern suggested that clinicians allowed patients to steer the dialogue to topics and issues initiated by the patient and gave patients enough time to hear them out. This interpretation is in keeping with the affective ratings, which indicated that during visits with patients in the intervention, clinicians sounded more engaged and sympathetic and no more hurried than during visits with control group patients, even though the visits were longer.
Beyond differences in specific categories of talk, the medical dialogue of visits involving a patient from the intervention group was characterized as significantly more patient centered than visits involving control group patients. Although patient-centered communication may be defined in multiple ways, there is broad consensus affirming its value as a conceptual marker of high-quality care and professionalism (16,17).
The relevance of patient centeredness to mental health care may be especially critical in fostering a therapeutic alliance. As noted by several investigators, patient centeredness and therapeutic alliance share conceptual ground (18,19). As reported by Wissow and colleagues (19), there is a significant and positive relationship between the RIAS-based patient-centeredness measure and therapeutic alliance as measured by the Vanderbilt Therapeutic Alliance Scale. This suggests that the activation intervention not only enhanced communication generally but may have created a milieu that fostered a strong affective bond and agreement on goals and processes of treatment.
The study's significant limitations included its small sample and limited generalizability. Fewer than one-third of clinicians at the study sites participated by referring at least one patient, a finding that suggests there was a sizable self-selection bias. Consequently the results were likely to represent clinicians and patients who were most interested in communication and patient activation. As in all observational studies, the presence of a device to record conversation may have influenced communication by inspiring best or most conscientious practice. It is unlikely, however, that best practice was systematically interpreted by therapists in a way that would jeopardize the validity of findings. The issue of performance bias in response to tape recording has been addressed by several studies (20–22). All have found that the effect is minimal.
Patients were the unit of randomization, and so providers conducted visits with patients in both the intervention and the control groups, and no attempt was made to blind clinicians to the patient's group. Despite the advantages and strengths of this design, it is possible that some providers may have diminished group differences through a form of compensatory communication with control group patients. For instance, providers may have attempted to elicit greater engagement in the dialogue by asking these patients more questions than they asked intervention patients.
Our findings affirm the potential utility and acceptability of a patient activation intervention designed to improve communication about treatment during visits with clinicians in community mental health settings. Web-based interactive tools can be made widely available and used in clinics and other settings. Patient feedback grounded in evidence-based treatment recommendations for schizophrenia encourages patients to discuss treatment concerns, questions, and options with their therapist. The feedback provides patients with information tailored to their current needs and symptoms. It allows them to discuss potential evidence-based treatment options to resolve current symptoms or seek future outcomes they value, for example, employment.
This research was supported by grant MH067189 from the National Institute of Mental Health.
Dr. Roter is the author of RIAS and holds the copyright for the system. Johns Hopkins University (JHU) also has rights to enhancements of the system. Neither Dr. Roter nor JHU collects royalties for use of the system in research conducted at the university or its medical institutions. Dr. Roter is a coowner of RIASWorks LLC, a company that provides coding services. The other authors report no competing interests. |
Improving Attic Thermal Performance
With metal roofing systems becoming more popular for new construction and retrofits, the Florida Solar Energy Center wanted to learn more about these roofing systems and their effect on thermal performance.
For homes in hot climates, roofing systems and attic thermal performance have a heavy impact on cooling energy use. Using insulation to control the heat flux from the attic into conditioned spaces through the ceiling is a known way to increase an attic’s thermal performance, but there are other problems that need to be dealt with.There is heat gain to the thermal distribution system when ducts are located in the attic.Also, leaky supply ducts can cause negative pressures within the house when the air handler is on.The negative pressures can cause hot air from the attic to be drawn into the conditioned space through gaps around recessed light fixtures or other bypasses, adding further to the home’s cooling load (see Figure 1).
With metal roofing systems becoming more popular for new construction and re-roofing, the Florida Solar Energy Center (FSEC) wanted to learn more about these roofing systems, and about their effect on attic temperature.And we were curious to find out how well metal roofing systems performed compared to traditional black-shingle roofing systems. During the summer of 2002,we performed testing on four finished and unfinished metal roofing systems and two roofing systems using traditional dark asphalt shingles. Our research shows that controlling attic air temperatures can be just as important as reducing ceiling heat flux during times of peak cooling loads (see “Not So Extreme Attic Example,” p. 14).We found that the metal roofing systems we tested generally perform well, although the performance of some of them degrades over time.
Side-by-Side Roof Testing
The experiments took place at our Flexible Roof Facility (FRF) in Cocoa, Florida.The FRF is a 24 ft x 48 ft frame building constructed with its long axis oriented east-west (see photo above).The roof and attic are partitioned to allow simultaneous testing of multiple roof configurations. The attic is sectioned into six 6-ft-wide test cells that are thermally separated by partitions.The partitions between the individual cells are well sealed to prevent air flow between cells and are insulated to R-20 using 3 inches of isocyanurate insulation.The gable roof has a 5/12 pitch and 3/4-inch plywood decking.With the exception of cell 2,R- 19 unsurfaced batt insulation is installed in a consistent fashion between the attic trusses in the test bays.The attic is separated from the conditioned interior by 1/2-inch gypsum board.The interior of the FRF is a single air-conditioned space.
The facility allows reconfiguration with different roofing products, and we’ve used it to examine different levels of ventilation and installation configurations for tile roofing.We’ve also compared reflective roofing, radiant barriers, and sealed attic construction. Our recent testing addressed several questions. (Note that 1:300 ventilation is 1 ft2 of attic ventilation, or net free vent area, per 300 ft2 of attic floor area.)
• What is the performance (ceiling flux and attic air temperatures) of a standard black asphalt shingle roof with 1:300 ventilation (the control cell)?
• How does the Galvalume metal roof compare in thermal performance with a galvanized metal roof?
• How does an ivory metal shingle roof perform compared with the roof with a lower solar reflectance that was installed the previous summer?
• How does an innovative double roof construction with an insulated roof deck, radiant barrier, and no attic ventilation perform compared with other types?
• How does a white standing-seam metal roof perform compared with an unfinished metal roof?
We used six different test configurations. (All the vented cells have soffit and ridge venting.)
Cell 1. A Galvalume 5-vee unfinished metal roof with a 1:300 vented attic (1st year tested).
Cell 2. Black asphalt shingles with a vented double roof deck, a radiant barrier, 6-inch foam insulation on the underside of the bottom roof deck, and an unvented attic (2nd year tested).
Cell 3. Solar-reflective ivory metal shingles with a 1:300 soffit and ridge ventilation (1st year tested).
Cell 4. A galvanized 5-vee unfinished metal roof with 1:300 ventilation (1st year tested).
Cell 5. Black asphalt shingles with a 1:300 soffit and ridge ventilation.This configuration has been installed for 15 years and is the control cell.
Cell 6. A white standing-seam metal roof with a 1:300 vented attic (7th year tested).
All the roofing materials were installed in a conventional manner according to the manufacturers’ specifications.
We tested samples of the new unexposed roofing materials to establish their solar reflectance and long-wave emittance (see Table).There is a large difference in the long-wave emittance of the two unfinished metal roofs.The emittance of Galvalume (0.28) is much lower than that of the painted metal roof (0.83), but the emittance of the unfinished galvanized roof is much lower still (0.04). Generally, low emissive surfaces reach higher temperatures than high emissive surfaces since low emissive surfaces absorb more solar radiation and do not readily give up collected heat back to the sky or to their surroundings.
Using precision thermocouples,we made a number of temperature measurements at the exterior surface of the roof and the underlayment; at the underside of the deck; in attic air at several heights within the attic; in the soffit inlet air and ridge vent exit air; at the top surface of the insulation; and at the interior ceiling of the conditioned space.The outside air temperature, insolation, humidity, and rainfall were also measured.All of the test cells were monitored from June 5 to September 30, 2002.
To find differences in space cooling, we evaluated the different roof systems to examine attic temperatures, heat flux through the ceiling into the conditioned space below, and the overall impacts of the thermal differences of the roofing system on required air conditioning.
Attic Air Temperatures
Controlling attic air temperatures is important, since the time when the attic becomes hottest—late afternoon—is almost always the time when the A/C runs the most and gains to the duct system are increased.The average summer day midattic air temperature profiles show the impact of the roofing options in reducing attic air temperatures and cooling energy use associated with attic duct heat gains and loads from unintended air leakage coming from the attic zone (see Figure 2).The statistics for the average, minimum, and maximum midattic air temperatures over the entire summer show that the sealed attic with the double roof provides the lowest overall mean attic temperatures (77.7ºF) and hence the lowest attic duct system heat gains and impact from return air leakage from the attic.
The next most productive roof combination in this regard is cell 6 with the vented white metal roof (81ºF). Cell 3, with the IR reflective metal shingle roof, had a very similar performance (82.2ºF). Next best in performance is cell 1, with the Galvalume metal roof and the vented attic at 83.6ºF.The low-emissivity galvanized metal roof (cell 4) was the least effective of the metal roof systems. Cell 4 attic temperature averaged 85.1ºF compared with the standard attic, which averaged 89.1ºF.
Maximum Attic Air Temperatures
We compared the average daily maximum midattic air temperature for each cell against the average daily maximum ambient air temperature along with the corresponding temperature difference for the full summer period.These results show the success of the various roofing options in controlling duct heat gains and loads from unintended air leakage under peak conditions for the period.
Note that cell 2, with the sealed attic and insulation on the underside of the roof decking cannot be directly compared with the other cells, because the other cells do not have roof deck insulation; instead, they have insulation on top of the ceiling.However, when we compared the 2002 summer results with the 1999 and 2000 cell 2 results (sealed attic without double roof deck and radiant barrier),we found that the average maximum midattic temperature difference from ambient was 4.7ºF lower for the double roof/radiant barrier combination than it was for the same sealed attic without the double roof.The maximum midattic temperature for the double roof deck and radiant barrier was 81.1ºF, or 7.1ºF lower than the averaged 1999 and 2000 results.The outdoor temperatures were quite similar during the two years in Florida’s predictably hot summers.
The high-reflectance ivory metal shingle (cell 3) provided the coolest attic of the test cells without roof deck insulation. The average maximum midattic temperature in this case was 93.3ºF, or 7.4ºF higher than ambient. In 2001, the brown, IR reflective shingle on the test cell had a maximum attic air temperature that was 10.6ºF higher than ambient. In 2000, the brown (not-high reflectance) metal shingle that was on the same cell had an average maximum
attic temperature 13.5ºF higher than ambient, while in 1999, a white highreflectance metal shingle on the same cell had an average maximum attic temperature 3.8ºF higher than ambient. Thus, the new ivory colored IR reflective shingle is better than all the tested metal tile products except the white standing-seam metal roof.
The white standing-seam metal roof (cell 6) was cleaned prior to the test for comparison with the pristine Galvalume and galvanized metal roofs.Comparison with the previous year shows the benefits of the cleaning and venting. In 2001, the average daily maximum attic air temperature above ambient was +14.4ºF, as compared to +7.8ºF in the summer of 2002.
Ceiling Heat Flux
The uninsulated ceiling of the double roof with sealed attic (cell 2) has a peak heat flux similar to that of the control (cell 5), although with a significant time lag of over 3 hours (see Figure 3).The mean heat flux for the double roof is 0.98 Btu/ft2 per hour, or 40% higher than the control.The double roof showed both the lowest mean and peak attic air temperatures of the group, and the highest ceiling heat flux.This seemingly contradictory result stems from the fact that the floor of the sealed attic (which forms the ceiling of the conditioned zone) is not insulated, so that the attic is unintentionally conditioned—reducing the attic temperatures by increasing heat transfer to the interior space.The absence of insulation on the attic floor produces the high heat fluxes to the interior.
The high reflectance ivory metal shingle roof (cell 3) has the lowest peak ceiling heat flux at 1.19 Btu/ft2 per hour. It also has a relatively low mean flux of 0.39 Btu/ft2 per hour, which is slightly higher than the mean flux for the white metal roof at 0.30 Btu/ft2 per hour.The vented white metal roof shows the lowest overall average heat flux and thus the lowest indicated ceiling influence on cooling for the overall period.The Galvalume roof (mean heat flux of 0.43 Btu/ft2 per hour) performs similarly to the IR reflective roof.The galvanized metal roof has poorer performance (mean = 0.53 Btu/ft2 per hour).
Overall Impact of Roofing System
The impact of roofing on cooling energy typically depends on three factors.These are the ceiling heat flux to the interior from the attic, the heat gain to the duct system located in the attic space, and the air unintentionally drawn from the attic into conditioned space. The heat flux through the ceiling impacts the interior temperature and hence the thermostat, which then calls for mechanical cooling.Thus, the heat flux affects cooling energy use at all hours, as well as the demand for air conditioning.
The other two influences—air leakage drawn from the attic into the conditioned space and heat gain to the duct system—usually occur only when the cooling system operates.Thus, the impact depends on the air conditioner run time in a particular time interval, and also on the leakiness of the duct, and on the amount of duct insulation. To obtain the average cooling system run time,we used a large set of residential cooling energy use data. These data come from 171 homes in the central Florida area, where the 15- minute air conditioner power was measured for over a year.
For each site, the maximum demand during summer was also recorded to determine the maximum cooling system power.Thus, it was possible to determine the diversified run time fraction by dividing the average air conditioner system power by its maximum demand.This calculation was made by averaging the air conditioner and air handler power for all sites and dividing by the average maximum summer demand, which was 3.96 kW.
The average cooling system run time is at its maximum (approximately 55%) at 4 pm (same as system diversity) and is at its minimum of (15%) at 6 am. It is important to note that these figures are based on an average summer day, as determined by evaluating all data from June 5 to September 30 inclusive.They do not represent extreme summer day conditions.
To estimate the impact of each roofing system on attic temperature, we assume a typical single-story home with 2,000 ft2 of conditioned floor area (see Figures 4 and 5).All of the alternative test cells do better than the control cell.The white metal roof with ventilation (cell 6) does best, followed by the high-reflectance metal shingle roof (cell 3).The Galvalume metal roof with a ventilated attic provides about a 30% reduction in heat gain.The galvanized roof with its significantly lower emissivity, provides only about a 20% heat reduction.The sealed attic with the double roof—cell 5—provides the lowest reduction in daily heat gain. This is primarily a result of the much greater measured heat flux across the uninsulated ceiling in this configuration.
As described earlier,we expect the unfinished galvanized steel roofing products to maintain their reflectance and emissivity properties less well over the long term than the Galvalume product.This is because the Galvalume’s aluminum-zinc alloy resists corrosion better.The preliminary data verify this expectation.We compared the maximum average daily attic air temperature for the summers of 2003 and 2002, looking at three types of metal roofing—white standing seam, galvanized, and Galvalume—and standard black asphalt shingle (see Figure 6).While the average maximum outdoor air temperature was 0.7ºF cooler,we found that each product showed some signs of weathering and increased solar absorptance, resulting in attic heating.
The average maximum attic air temperature under the standard black shingle roof showed no change; it was 116.7ºF in both years.However, the average maximum attic air temperature under the galvanized metal roof was 4.2ºF hotter in 2003 than in 2002.The average maximum attic air temperature under the Galvalume roof was 2.1ºF hotter in 2003, while the white metal roof showed an average increase of only 0.9ºF.Note that white metal remained the best choice with Galvalume next. This is consistent with anecdotal observation (Tennessee Williams’s Cat on a Hot Tin Roof ).After additional years of exposure,we expect that the Galvalume and galvanized will differ more widely in thermal performance.We expect that Galvalume will better maintain its performance, with most weathering occurring within the first year.This is because its surface properties are tailored to maintain its reflectance and weather resistance.The white metal roof was weathered in the sense that it accumulated dust and dirt, but there were no changes to the surface characteristics caused by oxidation.The white metal roof also showed the lowest degradation of the three metal roofs. Within the project, performance is being monitored for a third year in the same configuration to examine any further changes due to weathering.
Ranking the Roofs
Our test results from the summer of 2002 allowed us to compare the relative thermal performance of finished and unfinished metal roofing systems under typical Florida summer conditions (see Figure 7).The vented standing-seam white metal roof had the lowest total system heat gain of all the tested roofs, since its ceiling heat flux was much lower than that for roofs with the sealed attic construction. Its attic temperatures were also much lower than those for the standard black shingle control cell.The average daily maximum attic temperature was only about 94ºF. Cooling- related savings were on the order of 47% of roof-related heat gain.
The sealed attic double roof system (cell 2) provided the coolest attic space of all systems tested (average maximum daily midattic temperature was 81.1ºF), and therefore also the lowest estimated duct leakage and duct conduction heat gains.However, it also had the highest ceiling heat flux of all systems tested, due to the fact that the ceiling was uninsulated.This reduced its improvement over the standard black shingle roof in the control cell to a modest 7% savings in roof-related cooling energy. Note also that since this double roof configuration provided significantly cooler attic temperatures than the standard sealed attic tested during the previous two summers, higher total heat gains should be anticipated from standard sealed attics. Of course, it would be possible to combine both technologies—a cool roof and sealed attic construction—to produce even better results than any shown here.This suggests an area for future research.
A major objective of the testing was to evaluate popular unfinished metal roofing systems.We tested an unfinished Galvalume 5-vee metal roof with attic ventilation as well as a galvanized 5-vee metal roof in an identical configuration.The galvanized roof has a high solar reflectance, but a much lower long-wave emittance than the Galvalume roof (0.04 versus 0.28 for the Galvalume), which we expected to hurt its performance.The monitoring bore out this expectation.The Galvalume metal roof ran cooler than the galvanized system and produced less roof-related heat gain.The Galvalume roof provided a 32% reduction in roof and attic-related heat gain over the summer as compared with a 22% reduction for the galvanized roof. Moreover, as galvanized roofs are known to lose their solar reflectance rapidly over time as the zinc surface oxidizes,we expect to see a further decrease in performance in future seasons of testing.Although white metal performs best, the Galvalume metal roofing surface is a good second choice for cooling related climates, and does nearly as well as the IR selective ivory metal shingles.
At an average maximum midattic temperature of 93.3ºF (23.4ºF lower than the black shingle control cell), the high reflectance ivory metal shingle roof (cell 3) provided the coolest peak attic temperature of all the cells without a double roof deck.While its reflectance was somewhat lower than that of the white metal roof, the air space under the metal shingles provides additional effective insulation. Both of these characteristics probably come into play to help the high-reflectance roof achieve lower peak attic temperatures, while the additional insulating effect explains its slightly higher nighttime attic temperatures.
We also estimated the combined impact of ceiling heat flux, duct heat gain, and air being unintentionally drawn from the attic into conditioned space for the various roof systems. These estimates indicate that all of the tested roof systems yield lower heat gains during the summer cooling season than the control roof with black shingles.The rank order is shown in Figure 7, with the percentage reduction of roof/attic related heat gain and the approximate overall building cooling energy savings. It appears that nighttime attic temperature and reverse ceiling heat flux have a significant impact on total daily heat gain, and with greater benefit to constructions that produce lower evening attic temperatures. Since the roof/attic ceiling heat flux, duct heat transfer, and duct leakage probably comprise about one-third of the total home cooling load, the above values can be modified to approximate the overall impact.
The rank order of the reductions is consistent with the whole-house roof testing that we completed for FPL in Fort Myers, which showed that white metal roofing brings about reductions on the order of 20% of space cooling. However, these results represent the first time that popular unfinished metal roofs have been comparatively evaluated.
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|Pawnee Nation of Oklahoma tribal flag|
|Regions with significant populations|
|United States ( Oklahoma)|
|Related ethnic groups|
Pawnee people (also Paneassa, Pari, Pariki) are a Caddoan-speaking Native American tribe. They are federally recognized as the Pawnee Nation of Oklahoma and have four confederated bands: the Chaui, Kitkehakhi, Pitahawirata, and Skidi.
Historically, the Pawnee lived along outlying tributaries of the Missouri River: the Platte, Loup and Republican rivers in present-day Nebraska and in northern Kansas. They lived in permanent earth lodge villages where they farmed. They left the villages on seasonal buffalo hunts, using tipis while traveling.
In the early 19th century, the Pawnee numbered about 10,000 people and were one of the largest and most powerful tribes on the Great Plains. They had escaped some of the depredations of exposure to Eurasian infectious diseases impacting other Indian groups. By 1859, their numbers were reduced to about 1,400; however, by 1874 they were back up to 2,000. Still subject to encroachment by the Lakota and European Americans, finally most accepted relocation to a reservation in Indian Territory. This is where most of the enrolled members of the nation live today. Their autonym is Chahiksichahiks, meaning "men of men".
There are approximately 3,240 enrolled Pawnee, with 1,791 living in Oklahoma. Their tribal headquarters is in Pawnee and their tribal jurisdictional area is in parts of Noble, Payne, and Pawnee Counties. Their current elected president is Marshall R. Gover. Enrollment into the tribe requires a minimum 1/8th blood quantum.
Current Pawnee Business Council: Marshall R. Gover, President; Bruce Pratt, Vice President; Misty M. Nuttle, Treasurer; Phammie N. Littlesun, Secretary; Richard Tilden, Council Seat #1; Karla Knifechief, Council Seat #2; Adrian Spotted Horsechief, Council Seat #3; Liana Chapman Teter, Council Seat #4.
The new Council members were voted in by the people, elections are held every two years the first Saturday in May.
Economic development
They issue their own tribal vehicle tags, operate their housing authority, and maintain two casinos, three smoke shops, two fuel stations, and one truck stop. Their estimated economic impact for 2010 was $10.5 million.
Traditional culture
The Pawnee were divided into two large groupings—the Skidi living in the north and the South Bands (which were further divided into several villages). While the Skidi were the most populous group of Pawnee, the Chaui of the South Bands were generally the political leading group, although each band was autonomous. As was typical of many Indian tribes, each band saw to its own. In response to pressures from the Spanish, French and Americans, as well as neighboring tribes, the Pawnee began to draw closer together.
- Chaui, Chawi, or Tsawi (‘People in the Middle’, also called "Grand Pawnee")
- Kithehaki or Kitkehaxki (‘Little Muddy Bottom Village’, often called "Republican Pawnee")
- Pitahaureat or Pitahawirata (‘People Downstream’, ‘Man-Going-East’, derived from Pita - ‘Man’ and Rata - ‘screaming’, the French called them "Tapage Pawnee" - ‘Screaming, Howling Pawnee’, later the Americans "Noisy Pawnee")
- Pitahaureat, Pitahawirata, (Pitahaureat proper, leading group)
- Kawarakis (derived from the Arikara language Kawarusha - ‘Horse’ and Pawnee language Kish - ‘People’, some Pawnee argued that the Kawarakis spoke like the Arikara living to the north, so perhaps they belonged to the refugees (1794–1795) from Lakota aggression, who joined their Caddo kin living south)
- Turikaku (‘Center Village’)
- Kitkehaxpakuxtu (‘Old Village’ or ‘Old-Earth-Lodge-Village’)
- Tuhitspiat (‘Village-Stretching-Out-in-the-Bottomlands’)
- Tukitskita (‘Village-on-Branch-of-a-River’)
- Tuhawukasa (‘Village-across-a-Ridge’ or ‘Village-Stretching-across-a-Hill’)
- Arikararikutsu (‘Big-Antlered-Elk-Standing’)
- Arikarariki (‘Small-Antlered-Elk-Standing’)
- Tuhutsaku (‘Village-in-a-Ravine’)
- Tuwarakaku (‘Village-in-Thick-Timber’)
- Akapaxtsawa (‘Buffalo-Skull-Painted-on-Tipi’)
- Tskisarikus (‘Fish-Hawk’)
- Tstikskaatit (‘Black-Ear-of-Corn,’ i.e.‘Corn-black’)
- Turawiu (was only part of a village)
- Pahukstatu (‘Pumpkin-Vine’, did not join the Skidi and remained politically independent, but in general were counted as Skidi)
- Tskirirara (‘Wolf-in-Water’, although the Skidi-Federation got its name from them, they remained politically independent, but were counted within the Pawnee as Skidi)
- Panismaha (also Panimaha, by the 1770s this group of the Skidi had broken off and moved towards Texas, where they allied with the Taovayas, the Tonkawa, Yojuanes and other Texas tribes, was referred to as the Panimaha or Panismaha)
The Pawnee had a sedentary lifestyle combining village life and seasonal hunting, which had long been established on the Plains. Archeology studies of ancient sites have demonstrated the people lived in this pattern for nearly 700 years, since about 1250 CE.
The Pawnee generally settled close to the rivers and placed their lodges on the higher banks. They built earth lodges that by historical times tended to be oval in shape; at earlier stages, they were rectangular. They constructed the frame, made of 10-15 posts set some ten feet apart, which outlined the central room of the lodge. Lodge size varied based on the number of poles placed in the center of the structure. Most lodges had 4, 8 or 12 center poles. A common feature in Pawnee lodges were four painted poles, which represented the four cardinal directions and the four major star gods (not to be confused with the Creator.) A second outer ring of poles outlined the outer circumference of the lodge. Horizontal beams linked the posts together.
The frame was covered first with smaller poles, tied with willow withes. The structure was covered with thatch, then earth. A hole left in the center of the covering served as a combined chimney/smoke hole and skylight. The door of each lodge was placed to the east and the rising sun. A long, low passageway, which helped keep out outside weather, led to an entry room that had an interior buffalo-skin door on a hinge. It could be closed at night and wedged shut. Opposite the door, on the west side of the central room, a buffalo skull with horns was displayed. This was considered great medicine.
Mats were hung on the perimeter of the main room to shield small rooms in the outer ring, which served as sleeping and private spaces. The lodge was semi-subterranean, as the Pawnee recessed the base by digging it approximately three feet below ground level. This insulated the interior from extreme temperatures. Lodges were strong enough to support adults, who routinely sat on them, and the children who played on the top of the structures. (See photo above.)
As many as 30-50 people might live in each lodge, and they were usually of related families. A village could consist of as many as 300-500 people and 10-15 households. Each lodge was divided in two (the north and south), and each section had a head who oversaw the daily business. Each section was further subdivided into three duplicate areas, with tasks and responsibilities related to the age of women and girls, as described below. The membership of the lodge was quite flexible.
The tribe went on buffalo hunts in summer and winter. Upon their return, the inhabitants of a lodge would often move into another lodge, although they generally remained within the village. Men's lives were more transient than those of women. They had obligations of support for the wife (and family they married into), but could always go back to their mother and sisters for a night or two of attention. When young couples married, they lived with the woman's family.
Political structure
The Pawnee are a matrilineal people. Ancestral descent is traced through the mother, and, traditionally, a young couple moved into the bride's parents' lodge. People work together in collaborative ways, marked by both independence and cooperation, without coercion. Both women and men are active in political life, with independent decision-making responsibilities.
Within the lodge, each north-south section included roles for the three classes of women:
- Mature women (usually married and mothers), who did most of the labor;
- Young single women, just learning their responsibilities; and
- Older women, who looked after the young children.
Amongst the collection of lodges, the political designations for men were essentially between:
- the Warrior Clique; and
- the Hunting Clique.
Women tended to be responsible for decisions about resource allocation, trade, and inter-lodge social negotiations. Men were responsible for decisions which pertained to hunting, war, and spiritual/health issues.
Women tended to remain within a single lodge, while men would typically move between lodges. They took multiple sexual partners in serially monogamous relationships.
The Pawnee women were "skilled horticulturalists" and cooks, cultivating and processing ten varieties of corn, seven of pumpkins and squashes, and eight of beans. They planted their crops along the fertile river bottomlands. These crops provided a wide variety of nutrients and complemented each other in making whole proteins... In addition to varieties of flint corn and flour corn for consumption, the women planted an archaic breed which they called "Wonderful" or "Holy Corn", specifically to be included in the sacred bundles.
The holy corn was cultivated and harvested to replace corn in the winter and summer sacred bundles. Seeds were taken from sacred bundles for the spring planting ritual. The cycle of corn determined the annual agricultural cycle, as it was the first to be planted and first to be harvested (with accompanying ceremonies involving priests and men of the tribe as well.)
In keeping with their cosmology, the Pawnee classified the varieties of corn by color: black, spotted, white, yellow and red (which, excluding spotted, related to the colors associated with the four semi-cardinal directions). The women kept the different strains pure as they cultivated the corn. While important in agriculture, squash and beans were not given the same theological meaning as corn.
After they obtained horses, the Pawnee adapted their culture and expanded their buffalo hunting seasons. With horses providing a greater range, the people traveled in both summer and winter westward to the Great Plains for buffalo hunting. They often traveled 500 miles or more in a season. In summer the march began at dawn or before, but usually did not last the entire day.
Once buffalo were located, hunting did not begin until the medicine men of the tribe considered the time propitious. Then the hunt began by the men advancing together toward the buffalo, but no one could kill any buffalo until the warriors of the tribe gave the signal. Anyone who broke ranks was severely beaten. During the chase, the hunters guided their ponies with their knees and wielded bows and arrows. They could incapacitate buffalo with a single arrow shot into the flank between the lower ribs and the hip. The animal would soon lie down and perhaps bleed out, or the hunters would finish it off. An individual hunter might shoot as many as five buffalo in this way before backtracking and finishing them off. They preferred to kill cows and young bulls, as the taste of older bulls was disagreeable.
After successful kills, the women processed the bison meat and skin: the flesh was sliced into strips and dried on poles over slow fires and stored. Prepared in this way, it was usable for several years. Although the Pawnee preferred buffalo, they also hunted other game, including elk, bear, panther, and skunk, for meat and skins. The skins were used for clothing and accessories, storage bags, foot coverings, fastening ropes and ties, etc.
The people returned to their villages to harvest crops when the corn was ripe in late summer, or in the spring when the grass became green and they could plant a new cycle of crops. Summer hunts extended from late June to about the first of September; but might end early if hunting was successful. Sometimes the hunt was limited to what is now western Nebraska. Winter hunts were from late October until early April and were often to the southwest into what is now western Kansas.
Like many other Native American tribes, the Pawnee had a cosmology with elements of all of nature represented in it. They based many rituals in the four cardinal directions. Medicine men created sacred bundles which included materials, such as an ear of corn, with great symbolic value. These were used in many religious ceremonies to maintain the balance of nature and the Pawnee relationship with the gods and spirits. The Pawnee were not part of the Sun Dance tradition. In the 1890s, the people participated in the Ghost Dance movement.
The Pawnee believed that the Morning Star and Evening Star gave birth to the first Pawnee woman. The first Pawnee man was the offspring of the union of the Moon and the Sun. As they believed they were descendants of the stars, cosmology had a central role in daily and spiritual life. They planted their crops according to the position of the stars, which related to the appropriate time of season for planting. Like many tribal bands, they sacrificed maize and other crops to the stars.
The Morning Star ritual
The Skidi Pawnee practiced child sacrifice, specifically of captive girls, in the "Morning Star ritual". They continued this practice regularly through the 1810s and possibly after 1838, the last reported sacrifice. They believed the longstanding rite ensured the fertility of the soil and success of the crops, as well as renewal of all life in spring. The sacrifice was related to the belief that the first human being was a girl, born of the mating of the Morning Star, the male figure of light, and Evening Star, a female figure of darkness, in their creation story.
Typically, a warrior would dream of the Morning Star, usually in the autumn, which meant it was time to prepare for the various steps of the ritual. The visionary would consult with the Morning Star priest, who helped him prepare for his journey to find a sacrifice. With help from others, the warrior would capture a young unmarried girl from an enemy tribe. The Pawnee kept the girl and cared for her over the winter, taking her with them as they made their buffalo hunt. They arranged her sacrifice in the spring, in relation to the rising of the Morning Star. She was well treated and fed throughout this period.
When the morning star rose ringed with red, the priest knew it was the signal for the sacrifice. He directed the men to carry out the rest of the ritual, including the construction of a scaffold outside the village. It was made of sacred woods and leathers from different animals, each of which had important symbolism. It was erected over a pit with elements corresponding to the four cardinal directions. All the elements of the ritual related to symbolic meaning and belief, and were necessary for the renewal of life. The preparations took four days.
A procession of all the men, boys and male infants accompanied the girl out of the village to the scaffold. Together they awaited the morning star. When the star was due to rise, the girl was placed and tied on the scaffold. At the moment the star appeared above the horizon, the girl was killed with an arrow, then the priest cut the skin of her chest to bleed. She was quickly shot with arrows by all the participating men and boys to hasten her death. The girl was carried to the east and placed face down so her blood would soak into the earth, with appropriate prayers for the crops and life she would bring to all life on the prairie.
About 1820-1821, news of the sacrifices reached the East Coast; it caused a sensation among European Americans. Before this, US Indian agents had counseled Pawnee chiefs to suppress the practice, as they warned of how it would upset the American settlers, who were arriving in ever greater number. Knife Chief ransomed at least two captives before sacrifice. For any individual, it was extremely difficult to try to change a practice tied so closely to Pawnee belief in the annual renewal of life for the tribe. In June 1818, the Missouri Gazette of St. Louis contained the account of a sacrifice. The last known child sacrifice was of Haxti, a 14-year-old Oglala Lakota girl, on April 22, 1838.
Writing in the 1960s, the historian Gene Weltfish drew from earlier work of Wissler and Spinden to suggest that the sacrificial practice might have been transferred in the early 16th century from the Aztec of present-day Mexico. More recently, historians have disputed the proposed connection to Mesoamerican practice. They believe that the sacrifice ritual originated separately within ancient traditional Pawnee culture.
Francisco Vásquez de Coronado visited the neighboring Wichita in 1541 where he encountered a Pawnee chief from Harahey in Nebraska. Nothing much is mentioned of the Pawnee until the 17th and 18th centuries when successive incursions of Spanish, French and English settlers attempted to enlarge their possessions. The tribes tended to make alliances as and when it suited them. Different Pawnee subtribes could make treaties with warring European powers without disrupting the underlying unity; the Pawnee were masters at unity within diversity.
In the 18th century, the Pawnee were allied with the French, with whom they traded. They played an important role in halting Spanish expansion onto the Great Plains by decisively defeating the Villasur expedition in battle in 1720.
Until the 1830s, the Pawnee in what became United States territory were relatively isolated from interaction with Europeans and escaped some of the losses due to introduction of Eurasian infectious diseases, such as measles, smallpox, and cholera, to which Native Americans had no immunity. In the 19th century, however, they were pressed by Siouan groups encroaching from the east, who also brought disease. Epidemics of smallpox and cholera, and endemic warfare with the Sioux and Cheyenne drastically reduced the numbers of Pawnee. From an estimated population of 12,000 in the 1830s, they were reduced to 3,400 by 1859, when they were forcibly constrained to a reservation in modern day, Nance County, Nebraska.
In 1874, the Pawnee requested relocation to Indian Territory (Oklahoma), but the stress of the move, diseases and poor conditions on their reservation reduced their numbers even more. During this time, outlaws often smuggled whisky to the Pawnee. The teenaged female bandits Little Britches and Cattle Annie were imprisoned for this crime.
By 1900. the Pawnee population was recorded by the US Census was only 633. Since then the tribe has begun to recover in numbers.
The historian Marcel Trudel documented that close to 2,000 Pawnee (Panis in French) slaves lived in Canada until the abolition of slavery in the colony in 1833. The Indian slaves comprised close to half of the known slaves in French Canada (also called Lower Canada). Traditionally Native American and First Nations tribes sold captives from warfare as slaves to other tribes and to European traders.
In French Canada, Indian slaves were generally called Panis (anglicized to Pawnee), as most had been captured from the Pawnee tribe. Pawnee became synonymous with "Indian slave" in general use in Canada. As early as 1670, a historical reference was recorded to a Panis in Montreal.
A Pawnee tribal delegation visited President Thomas Jefferson. In 1806 Lieutenant Zebulon Pike, Major G. C. Sibley, Major S. H. Long, amongt others, began visiting the Pawnee villages. Under pressure from Siouan tribes and European-American settlers, the Pawnee ceded territory to the United States government in treaties in 1818, 1825, 1833, 1848, 1857, and 1892. In 1857, they settled on the Pawnee Reservation along the Loup River in present-day Nance County, Nebraska, but managed to keep their regular pattern of life. They were subjected to continual raids by Lakota from the north and west. On one such raid, a Sioux war party of over 1,000 warriors ambushed a Pawnee hunting party of 350 men, women and children. The Pawnee had gained permission to leave the reservation and hunt buffalo. About 70 Pawnee were killed in this attack, which occurred in a canyon in present-day Hitchcock County. The site is known as Massacre Canyon. Because of the ongoing hostilities with the Sioux and encroachment from American settlers to the south and east, the Pawnee decided to leave their Nebraska reservation in the 1870s and settle on a new reservation in Indian Territory, located in what is today Oklahoma.
Therefore warriors enlisted as Pawnee Scouts in the latter half of the 19th century in the United States Army. Like other groups of Indian scouts, Pawnee warriors were recruited in large numbers to fight on the Northern and Southern Plains in various conflicts against hostile native Americans. Because the Pawnee people were old enemies of the Sioux, Cheyenne, Arapaho, Comanche and Kiowa tribes, instead of fighting against Westward expansion, they served with the army for fourteen years between 1864 and 1877, earning a reputation as being a well trained unit, especially in tracking and reconnaissance. The Pawnee Scouts took part with distinction in the Battle of the Tongue River during the Powder River Expedition (1865) against Lakota, Cheyenne and Arapaho and in the Battle of Summit Springs. They also fought in the Great Sioux War of 1876. On the Southern Plains they fought against their old enemies, the Comanches and Kiowa, in the Comanche Campaign.
In 1875 most members of the nation moved to Indian Territory, a large area reserved to receive tribes displaced from the east and elsewhere. Many Pawnee men joined the United States Cavalry as scouts rather than face the ignominy of reservation life. The warriors resisted the loss of their freedom and culture by adapting to reservations. On November 23, 1892, the Pawnee in Oklahoma signed an agreement with the Cherokee Commission to accept individual allotments.
Recent history
The US government dismantled the Pawnee tribal government and civic institutions in 1906. The tribe reorganized under the Oklahoma Indian Welfare Act of 1936 and established the Pawnee Business Council, the Nasharo (Chiefs) Council, and a tribal constitution, bylaws, and charter.
In the 1960s, the government settled a suit by the Pawnee Nation regarding their compensation for lands ceded to the US government in the 19th century. By an out-of-court settlement in 1964, the Pawnee Nation was awarded $7,316,097 for land ceded to the US and undervalued by the federal government in the previous century.
Bills such as the Indian Self-Determination and Education Assistance Act of 1975 have helped address the mistakes of the past. The Pawnee Nation has regained some of its self-government, culture and pride. The Pawnee continue to practice cultural traditions, meeting twice a year for the inter-tribal gathering with their kinsmen the Wichita Indians. They have an annual four-day Pawnee Homecoming for Pawnee veterans in July. Many Pawnee also return to their traditional lands to visit relatives and take part in scheduled powwows.
Notable Pawnee
- Acee Blue Eagle, artist and educator
- Big Spotted Horse, 19th-century warrior and raider
- Larry Echo Hawk, Bureau of Indian Affairs Director. He was elected Attorney General of Idaho (1991–1995). He is a member of the First Quorum of the Seventy of The Church of Jesus Christ of Latter-day Saints.
- Walter Echo-Hawk, worked on major legislation, such as the Native American Graves Protection and Repatriation Act
- John EchoHawk, lawyer and founder of the Native American Rights Fund
- Kevin Gover, director of the National Museum of the American Indian.
- Moses J. "Chief" Yellow Horse, Major League baseball player
- Old Lady Grieves The Enemy, 19th-century woman warrior
- Petalesharo, a Skidi Pawnee chief who in 1817 rescued an Ietan Comanche girl from Pawnee ritual human sacrifice.
- Anna Lee Walters (b. 1946), Otoe-Missouria-Pawnee author and educator
- Wicked Chief, visited President James Monroe in 1822 with a delegation of Indian dignitaries.
See also
- 2011 Oklahoma Indian Nations Pocket Pictorial Directory. Oklahoma Indian Affairs Commission. 2011: 27. Retrieved 24 Jan 2012.
- Parks, Douglas R. "Pawnee." Oklahoma Historical Society's Enyclopedia of Oklahoma History and Culture. (retrieved 14 Sept 2011)
- "Preamble." Constitution of the Pawnee Nation of Oklahoma. Revised 14 June 2008.
- "Pawnee" Encyclopedia of Oklahoma History and Culture http://digital.library.okstate.edu/encyclopedia/entries/P/PA022.html, accessed 10 Dec 2012
- Hyde 22
- Weltfish 5
- Hyde 361
- Weltfish 463
- Weltfish 4–8
- Carleton, James Henry (1983). The Prairie Logbooks. Lincoln, Nebraska: University of Nebraska Press. pp. 66–68. ISBN 0-8032-6314-7.
- Weltfish 119–122
- W.P. Clark, "Hunt", The Indian Sign Language, Lincoln, NE: University of Nebraska Press (1982, first published 1885), trade paperback, 444 pages, ISBN 0-8032-6309-0
- Weltfish 106-118
- Weltfish 117
- Philip Duke, "THE MORNING STAR CEREMONY OF THE SKIRI PAWNEE AS DESCRIBED BY ALFRED C. HADDON"], The Plains Anthropologist, Vol. 34, No. 125 (August 1989), pp. 193-203
- Hyde 85–336
- "History of Nance County, Nebraska, NEGenWeb Project". Usgennet.org.
- "Cattle Annie & Little Britches, taken from Lee Paul [http://www.theoutlaws.com]". ranchdivaoutfitters.com. Retrieved December 27, 2012.
- Weltfish 3-4
- Carter Godwin Woodson, "The Slave in Canada", The Journal of Negro History, Vol. 5, July 1920, No. 3, pp. 263-264. Retrieved 17 November 2009
- Deloria Jr., Vine J; DeMaille, Raymond J (1999). Documents of American Indian Diplomacy Treaties, Agreements, and Conventions, 1775-1979. University of Oklahoma Press. pp. 361–363. ISBN 978-0-8061-3118-4.
- 78 Stat. 585 (1964); Wishart, David J., 1985. "The Pawnee Claims Case, 1947–64," Irredeemable America: The Indians' Estate and Land Claims, ed. I. Sutton (Albuquerque: University of New Mexico Press): 157–186.
- "Nominee Named for Indian Affairs", Associated Press, New York Times, 10 April 2009
- Hyde, George E. The Pawnee Indians. Norman: University of Oklahoma Press, 1974. ISBN 0-8061-2094-0.
- Weltfish, Gene. The Lost Universe, Pawnee Life and Culture. Lincoln: University of Nebraska Press, 1977. ISBN 0-8032-5871-2.
Further reading
- Robert O. Lagace, "Pawnee: Culture summary", Ethnographic Atlas, University of Kent, Canterbury.
- "Pawnee", Encyclopedia of North American Indians, New York: Houghton Mifflin
- J. S. Clark, "A Pawnee Buffalo Hunt", Oklahoma Chronicles, Volume 20, No. 4, December 1942, Oklahoma State Historical Society
|Wikimedia Commons has media related to: Pawnee|
- Pawnee Nation, official website
- Pawnee Indian Tribe, Access Genealogy
- Pawnee Indian History in Kansas
- Pawnee Indian Village Museum, A museum featuring the excavated floor of a large 1820s Pawnee earth lodge and associated artifacts, Kansas State Historical Society
- Non-invasive imagery of a Pawnee archaeological site, mapping the archaeological remains of a late 18th- and early 19th-century Pawnee village site located on the Republican River in north central Kansas.
- Inventory of the Gene Weltfish Pawnee Field Notes, 1935 at the Newberry Library |
There are two ways we can entrust our
money to a bank:
These are two very different types of transactions. They are quite
different in nature and have a different legal basis.
The Demand Deposit
When I get paid, my salary goes into my cheque account. I do not want
to lend my salary to the bank, because I intend to spend most of my pay
in the fortnight before I next get paid. I want to be able to spend that
money whenever I choose, for whatever I want to buy, so I choose a
When my salary goes into the bank, I have not transferred the
ownership of that money to the bank. The money is mine. It does not
belong to the bank.
The current accounting process is that the bank records the cash it
has received as an asset and records it responsibility to me as an
liability. This is wrong. The cash does not belong to the bank. It
belongs to me. The cash will be an asset on my balance sheet, so it
should not be an asset on the bank’s accounts at the same time. An
asset cannot have two owners.
Consider a parallel example. If I am going out of town for a while, I
may engage a warehousing company to store my dining suite until I
return. The warehouse will charge a fee for providing this service. When
my dining suite goes into the warehouse, its ownership does not change.
The dining suite still belongs to me. The warehouse owner cannot do what
he likes with the table. He cannot bring it out and use it when he has
guests for dinner. He cannot dance on the table top or use it for
playing table tennis. The warehouse owner cannot decide how the table
will be used, because he has no ownership rights to it. He has a duty to
care for my dining table in the way specified in the contract.
If I decide not to return, I can write to the warehouse and ask that
the table be delivered to my daughter and the chairs to one of my
friends. The warehouse owner will do this, provided I pay the cost of
transport. He cannot refuse to carry out my request, because he has
relatives staying and is using the table. If this happened, I would
accuse him of misappropriating my dinning suite. If he has moved it to
his own home, I could accuse him of theft. Everyone would understand
that he has done something immoral.
If the word got out about what he had done, his warehouse would soon
be empty, because people would stop trusting him. The service that he
offers is skill at caring well for things that belong to other people.
This service only has value to customers, if they can trust him to
provide the care that he has promised. He is really selling trust, so if
he proves to be untrustworthy, his service has no value and people will
be unwilling to pay for it.
The warehouse owner does not record the things stored in his
warehouse on his balance sheet. The only asset on his balance will be
the warehouse that he owns. He does not include the contents of the
warehouse, because he does not own them. They are not his assets. The
only way that they will appear on his balance sheet is through a
contingent liability for inadvertent damage that might be done to
something that is in his care.
When I put my salary in the bank, I am really just putting it in a
warehouse for safekeeping. This was obvious in the days when gold coins
circulated. Keeping the coins at home would be too risky, if I could not
afford a safe. I would be very vulnerable to being robbed. Putting my
gold coins in a bank for safe keeping would make good sense. The bank’s
service would be even better, if it would act on my instructions to make
payments to other people on my behalf when required. The bank would
simply transfer the ownership of the right number of coins to the person
to whom I was making a payment.The bank is just providing a warehouse
role for the gold coins. The ownership of the coins does not shift to
the bank, when I deposit them. The coins are still mine. They only
change ownership, when I instruct the bank to make a payment to someone.
At that point the ownership of some of the coins transfers to the person
paid. The bank never owns the coins. Therefore the gold coins should
never be recorded as an asset on the banks balance sheet. All that
should appear on the banks balance sheet is any contingent liability for
coins that are lost or stolen.
Bank of Amsterdam
The bank of Amsterdam operated in this way for more than 150 years
from 1609 to 1779. Very few banks have remained sound for this length of
time. The integrity of this bank was noted by David Hume and Adam Smith.
The latter made the following comments.
The Bank of Amsterdam professes to lend out no part of what is
deposited with it, but, for every guilder for which it gives credit in
its books, to keep in its repositories the value of a guilder either in
money or bullion. That it keeps in its repositories all the money or
bullion for which there are receipts in force, for which it is at all
times liable to be called upon, and which, in reality, is continually
going from it and returning to it again, cannot well be doubted...... At
Amsterdam, however, no point of faith is better established than that
for every guilder, circulated as bank money, there is a correspondent
guilder in gold or silver to be found in the treasure of the bank.....
The bank is under the direction of the four reigning burgomasters who
are changed every year. Each new set of burgomasters visits the
treasure, compares it with the books, receives it upon oath, and
delivers it over, with the same awful solemnity, to the set which
succeeds; and in that sober and religious country oaths are not yet
In 1672, when the French king was at Utrecht, the Bank of Amsterdam
paid so readily as left no doubt of the fidelity with which it had
observed its engagements. Some of the pieces which were then brought
from its repositories appeared to have been scorched with the fire which
happened in the town-house soon after the bank was established. Those
pieces, therefore, must have lain there from that time.......
The bank cannot be debtor to two persons for the same thing
Inquiry into the Nature And Causes of the Wealth of Nations, 1776).
The last sentence is the key to the longevity of the Bank of
Unfortunately, during the 1780s, the bank began to violate this
principle, when the city of Amsterdam demanded that surplus deposits be
loaned to the city. This marked the end of the bank's success.
What applies to gold coins applies to all forms of money. I get my
salary paid into a bank account because I do not want to carry around a
whole lot of notes and coins. I put it there for safekeeping. I also use
the bank because it provides an easy way of making payments to other
people, by cheque or electronic transfer.
I am very clear about one thing. My money belongs to me, even when it
is in the bank. I want to be able to spend it whenever I want. My cash
in the bank is my asset. It does not belong to the bank. Therefore the
bank should not record my cash on its balance sheet as an asset, even if
also records a liability to me. My cash does not belong to the bank.
This is the heart of the problem with the modern banking system.
Banks claim ownership of the cash that has been deposited by their
customers. They record this cash as an asset on their balance sheet.
They treat the cash as if they owned it. This is problematic because the
cash now has two owners. I think that I own it, while the bank acts as
if it owns it. (For a full analysis of the nature of money, see the
Having on asset with two owners might be fine for a while, if the
real owner does not want to use the asset immediately. However,
eventually problems will arise. If the owners of money in the bank want
to withdraw it and the bank has done something else with the money, the
conflict is obvious. If too many people want to withdraw at the same
time, the problem compounds. A bank run can occur, and the bank might
end up defaulting on its obligations.
This problem does not arise with a warehouse. If all the people with
stuff stored in a warehouse decided to take it out on the same, this
would not matter. The warehouse owner would be very busy handing out
stuff and he might be worried about his future income, but every person
would get back what they owned. There is no reason why a bank should be
The solution to this is quite clear. If the warehouse owner claimed
ownership of the stuff stored in his warehouse, he would be accused of
misappropriation or theft. If a bank claims ownership of money entrusted
to its care, the same applies. It has appropriated something that does
not belong to it. It has stolen money that it does not own.
The Bible is clear that a thing cannot have two owners. If two people
claim the same thing, the case should be resolved by judges.
In all cases of illegal possession of an ox, a donkey, a sheep, a
garment, or any other lost property about which somebody says, 'This is
mine,' both parties are to bring their cases before the judges. The one
whom the judges declare guilty must pay back double to his neighbor (Ex
If I say of my demand deposit ,'This is mine' and the bank is also
saying, 'This is mine,' something is wrong. This is an issue that should
be resolved by judges. If they find that the bank has claimed something
that does not belong to it, the bank pay back double to the depositor.
The warehouse owner will keep an inventory of everything that is
stored in his warehouse. He records the identity and contact details of
the owner of each item. He can even transfer the ownership to another
person, if instructed by to do so by the owner. However, this recording
system will be separate from his asset register.
Banks should really be doing the same thing. They should be keeping
an inventory of all the money being stored and the identity of the
owners. This should be separate from their financial accounts. The money
stored should not be able to creep onto the banks asset register.
Some defenders of the modern banking system have argued that a bank
is different from a warehouse, because money is fungible. A fungible
good is one that becomes mixed when it is stored with more of the same
good. An example is wheat. If my wheat is added to a silo containing
wheat owned by other people, I will not be able to get my particular
grains of wheat out. All that the silo owner is required to do is return
wheat of similar quality to what I put into a silo.
Oil is another example of a fungible good. Having a separate storage
tank for each owner of oil would not be very efficient. It is better to
store all the oil of a particular type in one tank. The operator of the
tank will always be able to return oil to the person who put some in
when he wants it, however he will not get exactly the same molecules of
oil that he put in. He will be quite happy as long as he gets oil of the
The defenders of modern banks argue that money is a fungible good
too. This is true. If you deposit gold coins in a bank, you may have an
attachment to a particular coin, because you like the face that is
printed on it, but most people will happy to get back any coin, as long
it contains gold of the same weight and purity. Likewise, I don’t
really care what bank notes I get when I with draw money from the bank,
as long as they are not scruffy. I do not care, if I don’t get back
the ones that I put in.
Money is a fungible good, but the argument still has a fatal flaw.
The defenders of modern money argue that when a fungible food is put
into a silo or tank, the ownership of the good passes to the owner of
the storage. They say that the same applies to demand deposits at the
bank. Because money is a fungible good, ownership passes to the bank.
This logic is not correct.
When oil is put in a silo, or oil into a tank, the ownership should
not pass to the owner unless he buys it from the people putting it in.
If I put wheat in a silo along with wheat that belongs to eight other
people, I still own some wheat. The difference is that whereas,
previously I owned my own wheat, I have now own a share of silo of
wheat. The wheat in the silo does not belong to the owner of the silo.
It is owned jointly by the other people who have put their wheat into
the silo. The owner of the silo is not entitled to take some of the
wheat for his own use, even if he puts some other wheat back, before the
owners demand it back.
Likewise, when several people put their oil into the same tank, the
ownership of the oil should not transfer to the owner of the tank.
Rather they have swapped ownership of a specific volume of oil for a
share in a larger volume of oil. The owner of the tank can only claim
ownership of oil if he has actually purchased the oil. He has no right
to take the oil for his own use, even if he puts some oil volume back
into the tank later. The owner of the tank or the silo has a duty of
care to those who are paying a fee for the storage. That does not give
him a right use what is being stored.
The same applies to money deposited in a bank. Although the money is
fungible, ownership does not transfer to the bank. Rather the depositors
own a share of all the money in the bank. This is true regardless of the
form of money. If gold is deposited in the bank, the depositor changes
ownership of a particular piece of gold, for a defined share of all the
gold in the bank. He can always get his gold back, because the amount he
has put in has been added to the total amount in the bank. Each
depositor’s share of the total amount of gold is equivalent to the
amount that they put in.The same applies if the money is notes and coins
or electronic money. All the money in the bank is jointly owned by the
depositors. Each one owns a share of the total, which is equivalent to
the amount they deposited. If the money has depreciated in value, the
loss is shared by all depositors. The important point is that none of
the deposited money is owned by the bank.
This is a biblical principle.
If a man gives his neighbor silver or goods for safekeeping and they
are stolen from the neighbor's house, the thief, if he is caught, must
pay back double. But if the thief is not found, the owner of the house
must appear before the judges to determine whether he has laid his hands
on the other man's property…. The one whom the judges declare guilty
must pay back double to his neighbor. (Ex 22:7-9).
When someone takes the
goods of another for safekeeping and it goes missing, he is accountable
for the loss. If the thief is found, the thief must make restitution. If
not, the person caring for the property is accountable for the loss. He
must make restitution to the owner, because his neglect is the
equivalent of theft.
The other important thing to note is that the Bible refers to the
valuables presented for safekeeping as the “property” of the
depositor, even when they are in the house of the other person. This
confirms the principle that the ownership of property does not transfer
to person who takes it for safekeeping. The owner of the property
remains the owner, until the goods are actually sold.Applying this
principle to banking, the bank that treats money that has been deposited
for safekeeping as its own asset has misappropriated something that does
not belong to it. It has “laid its hands on the other man’s property”.
If it was taken before the judges, it would have to pay back double to
the owner. Paying back the amount that was deposited is not enough.
There are limits on the duty of care that must be provided. The
person providing safekeeping is not accountable for events beyond their
If a man gives a donkey, an ox, a sheep or any other animal to his
neighbor for safekeeping and it dies or is injured or is taken away
while no one is looking, the issue between them will be settled by the
taking of an oath before the LORD that the neighbor did not lay hands on
the other person's property. The owner is to accept this, and no
restitution is required. But if the animal was stolen from the neighbor,
he must make restitution to the owner. If it was torn to pieces by a
wild animal, he shall bring in the remains as evidence and he will not
be required to pay for the torn animal (Ex 22:10-13).
The principle remains the same. The owner is the owner. The neighbour
providing care is never the owner. If the animal is stolen, the
neighbour must make restitution to the owner. If the animal is killed by
wild animals, the neighbour does not have to make restitution, because
this event was beyond his control.
The same applies to a bank. If it claims money that has been
deposited as its own asset, it has committed theft. If the bank is
robbed, it must make restitution to the depositor. However, if the money
is destroyed by a fire or war, the bank is not liable for the loss,
because it was beyond the bank’s control. A bank must provide the best
care possible for money deposited, but it is not accountable for events
beyond its control.
Banks that operate according to these biblical principles will not be
able to make money from accepting money on deposit for safekeeping.
Therefore, it will be quite appropriate for them to charge a fee for the
services that they provide. Depositors will look for banks that provide
the best service for the most reasonable fee. They will be able to
choose the level of service. Banks that provide better security and a
wider range of transactions will be able to charge more.
Banks should not be expected to pay interest on demand
Term Deposits or Loans
Sometimes when I receive my salary, I will not need to spend it all
in the next fortnight. I might build up the size of my deposits in the
bank to the level where I have a couple of thousand dollars that I do
not want to spend until some time in the future.
In this case, I might want to put the money on a term deposit, so I
can earn interest as a reward for not using it immediately. This is a
different type of transaction. I am agreeing to let the bank have the
use of the money for a specified time. I am allowing the bank to decide
how to lend the money. I expect them to do this in a way that will keep
my money safe, while earning a good return.
This is different from the demand deposit, so the analogy of the
warehouse does not apply. This transaction is more akin to someone who
leases a piece of machinery. If a person owns a tractor, but has no way
to make use of it, the tractor brings them no benefit. If they can find
a farmer who needs, a tractor, they can lease it to the farmer for use
on his farm. The farmer is better off because he has the use of a
tractor. The tractor owner is better off, because his tractor now earns
him some income.
The same applies with cash sitting in demand deposit for month after
month. If I do not want to spend it immediately, I am better to lease it
to someone else who can use it productively. They are better off because
they have capital that they did not have. I am better off, because they
will pay me for the use of the money. This is the nature of a term
deposit. I am making a loan to someone who can use it more efficiently
in return for interest. The interest is really just the rent or lease
for the use of the money.
Ownership and Risk
Some defenders of the modern financial system would argue that money
loaned to the bank belongs to the bank. This is not correct. The money
loaned still belongs to the person who saved it. The bank just has the
use of the money for a specified time.
This is obvious in the case of the tractor. When a farmer leases a
tractor, it does not become the property of the farmer. What he buys is
the right to use the tractor for a specific time. The farmer will not
list the tractor as an asset on his balance sheet. Rather he will record
the rent paid for leasing the tractor as an expense in his Profit and
Loss Account. Ownership of the tractor does not pass to the farmer
(except in the case of a lease to buy agreement).
The difference from the warehouse situation is that the owner can not
demand the tractor back whenever he chooses. He can only demand it back,
when the lease has expired. The other difference is that the farmer can
use the tractor how he likes, provided that he complies with the
conditions of the leasehold agreement. The lease may place limits on the
load he can pull and specify regular maintenance that must be done, but
provided the farmer complies with these conditions, he can use the
tractor how he pleases.
The principles for leasing productive equipment are confirmed in the
If a man borrows an animal from his neighbor and it is injured or
dies while the owner is not present, he must make restitution. But if
the owner is with the animal, the borrower will not have to pay. If the
animal was hired, the money paid for the hire covers the loss (Ex
This refers to a productive animal. Even when it has been lent to a
neighbour it still remains the property of the owner. The owner is till
Although the animal was leant for a fee or lease, the neighbour is
required to take reasonable care. However, if something unexpected goes
wrong, then he does not have to pay compensation. The lease fee will
have included something to cover the risk of the animal dying. The owner
must carry the risk of the animal dying. He would not lease the animal
unless the fee was sufficient to compensate for that risk.
Applying this to a tractor, the person leasing the tractor remains
the owner, even after it has been leased. The person leasing the
tractor, has control over its use, but does not have ownership. If the
engine unexpectedly blows up, the person leasing it is not responsible
(unless the lease specifies something different). The owner of the
tractor carries the risk of the tractor not performing as was expected.
The lease fee should compensate for that risk. If it did not, then the
owner would be foolish leasing the tractor.
These principles apply to a bank loans. When I loan money to a bank,
I still own the money. Therefore the bank should not record the money as
an asset in its balance sheet. The bank should pay interest for the use
of the money. Although the lender still owns the money, they cannot
demand the money back at any time. The lender must weight till the end
of the term for the money returned (unless he pays to break the
The bank has responsibility for using the money wisely, but this
cannot prevent problems occurring. Every transaction that involves the
future has uncertainty, so risk is unavoidable. The owner of the money
must bear this risk. The interest received should compensate the lender
for this risk. If it does not, then the lender would be better to leave
the money on demand deposit, so the banker carries the risk.
Money loaned to the bank is also fungible. When a person deposits
money in a bank, it will mix the money up with money deposited by other
people to make it more productive. This does not shift ownership of the
money to the bank. It just means that the depositor has shifted from
ownership of a specific amount of money, to ownership of a share of a
larger pool of money. The share of the money jointly owned will be equal
to amount of money deposited.
The important point is that ownership of the money does not transfer
to the bank. However the bank does have the right to use the money as it
chooses, provided it complies with the conditions of the deposit
agreement. Ownership is not transferred, but the right to control the
use of the money is transferred to the bank for a specific time.
The depositor owns the money, but the bank controls its use. It can
lend the money to someone else, but not for a longer period than the
loan agreement specifies. The loan agreement may also specify they types
of uses for which the money can be loaned by the bank. Provided these
conditions are met, the bank can decide how the money is used. If the
bank can get a higher interest rate than it agreed to pay to the
depositor, it can keep the difference to cover its expenses and provide
a profit to the bank owners. The bank cannot loan the money for a longer
term than the term agreed with the owners of the money, because it does
not have control over the money beyond that time. It only has control
for a specified time. When that period is complete, the bank loses this
control and must return the money to the owners.
Borrowing Short and Lending Long
Modern banks tend to borrow money for a short term and lend it for a
longer term. For example, most mortgages have a term of 15 to 20 years,
but much of the funding comes from term deposits with a term of less
than a year. This creates problems, because banks are lending money that
they do not own or control. If the bank only controls the money for 6
months lending it for a twenty year term is rather odd. We need to think
more clearly about this.
The depositor is in a strange situation. They have lent their money
to the bank for six months. The bank has taken this money and lent it to
someone for twenty years. It has taken this action knowing that it will
not get the money back from the borrower, when the term deposit is due
to be repaid. Unless the bank has other mortgages that will to be repaid
at that time, it will have to obtain the money to repay the depositor
from some other source.
A term deposit at a bank is supposed among the safest things that a
person can do with their money. This is not true. The reality is that
the bank is making a commitment to return the depositor’s money at the
time when the term is complete, without knowing where or how they will
obtain the money to fulfil this commitment. Most of this risk rests with
the bank. In normal times, they should be able to get the finance from
another source, but they cannot guarantee the future. Economic
conditions might have changed dramatically by the time the term deposit
matures. If the economy has deteriorated, the bank might have to offer a
much higher interest rate to get the finance they need.The security of a
term deposit depends on the ability of the bank to guess what will
happen in the future. A prudent bank might be able to manage the risks
involved, but if it guesses wrong about the future, it might face
significant losses. If the bank really gets thing wrong, and the losses
are greater than the banks capital, the depositors will bear the loss
and lose some of their money.
When a bank accepts a term deposit and lends the funds for a longer
term, it does not know where the funds will come from to repay the
depositor when the term is complete. It is committing to obtain money in
the future without knowing whether this will be possible or what the
price will be. It is making a commitment to return something that it
currently does not own.
A market where people buy and sell entities, which they do not own,
is called a futures market. Futures markets exist for a number of
commodities and for a variety of financial instruments. Such a market is
legitimate for people who want to speculate on or hedge against future
changes in price. All participants in the market understand that the
person who is promising to sell at a future date does not currently own
what they are promising to sell. There is a risk that when they try to
buy what they have agreed to sell; the price may have risen so high,
that they cannot afford to buy it. If that happens they might default on
their agreement. Everyone in the market understands this risk.
Futures markets are inherently unstable, as prices can fluctuate
rapidly and players default if unexpected events occur. This is fine for
investors who understand the nature of market and are prepared to take
on the risks involved. However, this type of arrangement will rarely be
appropriate for people making term deposits in a bank. They put their
savings in term deposits at a bank because they want their money to be
safe. They are generally willing to get accept lower interest rates in
return for higher security. They are not choosing to lend to an entity
that is speculating in a futures market. Depositors might be able to get
a higher rate of interest from a term deposit for a bank that borrows
short and lends long, bit if they understood the risk, most would prefer
a safer option.
Trust and Honesty
If bank depositors understood the risks involved, they would not have
so much confidence in the banks. If the farmer who borrowed the tractor
for six months lent it to a contractor for 10 years, the owner would get
a bit upset. He has lent his tractor to a farmer who does not know how
he will be able to return the tractor when the six months is up. The
farmer might be able to borrow or buy an equivalent tractor, but that
might quite costly. The owner of the tractor would become quite nervous
and would regret leasing the tractor to this farmer. He would probably
consider the farmer to be dishonest.
Depositors should view banks in the same way. Why should they trust a
bank that has taken money on deposit for six months, and lent it to
someone else for 10 years? Why should depositors trust the bank, when
the bank does not know how it will repay the money when the term of
their deposits are complete? Why are banks that behave in this way not
labelled dishonest? We should demand better behaviour from the banks
that care for our money. We should not trust organisations that behave
Most depositors would prefer to deposit their money with a bank that
only makes loans for a term that matches the term of their deposits.
Every loan issued by the bank would be matched by a deposit or group of
deposits with the same term. All loans would be for a fixed term and the
bank could only make loans with a term matching the terms deposits
already received by the bank. This is a much safer policy, as the bank
would know that borrowers will be repaying their loan when the term
deposits come to maturity. The bank would know where the money that
belongs to the depositors has gone and it would know when it will be
coming back. This would be a much safer way to operate, as much of the
risk is eliminated.
This policy of matching loans would have the effect of raising
interest rates on longer term deposits. This is reasonable, as the
longer the term of the loan, the greater is the risk of loss. Very long
term loans would probably disappear altogether. This might also be a
good thing. The Bible suggests that Christians should not take loans
with a term of more than seven years. The reason is that we do not know
what the future holds. Making commitments to do things in twenty or
thirty years time is very unwise, as we simply do not know what our
situation will be that far into future. Limiting loans to a maximum of
seven years would be a sensible policy for a bank.
Depositors need greater choice. A bank that only made loans for terms
that are matched by the terms of deposits they have received from
depositors would be much safer. A bank that is serious about providing
services to their customers should provide this service. People choose
bank term deposits, because they prefer security to high returns. I
expect they would prefer a bank that matches the terms of loans and with
the terms of deposits.
The interest rate for deposits with matched loans might be lower, but
the risk would be lower as well. If more people demanded this service,
some banks might start offering it. They would essentially become loan
brokers. They would be pooling deposits and matching these up with
borrowers wanting funds for the same term. They might charge a fee for
this service, or they might take add a margin on to the interest rate
paid to the depositors.
Risk of Default
Matching the term of the loan to the term of the deposits does not
eliminate all risk. The person who borrowed the money may abscond or
make bad business decisions. They might not be able to repay the money
they have borrowed when the term is complete. The bank is to should be
skilled in assessing the creditworthiness of borrowers and putting
appropriate security measures in place.
The more risky the loan, the higher the interest rate will have to
be. Depositors should be able to choose the level of risk they want to
take on, but I suspect that most depositors would specify that their
money only be lent to creditworthy borrowers.
The bank could also provide insurance against the borrower
defaulting. This does not eliminate the risk, but it spreads the cost
across all depositors, rather than leaving all the risk with the few
depositors affected by the bad loan. Most depositors would prefer a
slightly lower interest rate, if they knew that the cost of any default
would be spread across many depositors.
The person borrowing from a bank that lends long and borrows short
also faces uncertainty. They have agreed to a mortgage with at twenty
year term, with a bank that has only organised the finance for the first
six months of the loan. The bank presumes that it will continue to be
able obtain the money as it is needed. This is probably true, but the
bank does not know what rate of interest rate it will have to pay to
obtain the money in the future. That is why most banks will generally
not hold mortgage interest rates fixed for more then about two years.
They do not know how the interest rates will move in the future, so they
pass the risk on to the borrower.
The adjustable or floating rate mortgage is the solution to this
problem. The adjustable rate mortgage is really a series of short-term
loans. The bank is implicitly saying to the borrower, I can only loan
you this money for a couple of years. At then end of that term, I will
renew the loan, but I do not know what the interest rate will be. I will
always be able to borrow some money, but I cannot control the interest
rate. I will have to adjust the interest rate that you pay to allow for
this. So really an adjustable rate mortgage is really a series of short
term loans with different interest rates, but with the same security
remaining in place.
If the borrower is willing to take the risk of committing to a twenty
or thirty year mortgage, without know what the interest rate will be,
that is their business. However, taking on an unknown risk in this way
is not very wise. The Bible suggests that we should not make commitments
beyond five years. Making a commitment to make payments in thirty years
time is serious enough. The fact that you do not know what the amount
you have to pay will be makes the uncertainty even worse.
When Banks borrow short and lend long, depositors do not benefit.
Borrowers do not benefit either. The benefits all go to the bankers and
not their clients.
However, this policy is also risky for the bankers. The practice
works as long as the flow of withdrawals and deposits roughly match each
other. In a time of uncertainty or panic, a large number of people may
try to withdraw their deposits. As there is an advantage in getting in
first, there may be a run on the bank. If the bank cannot call up all
its loans it may run out of reserves. This could push the bank into
bankruptcy and many people would lose their money. Bank panics have been
common throughout the history of banking.
Modern banks are not very clear about the contracts they offer. Some
savings accounts pay interest, but have a variable term allowing the
money to be withdrawn on demand. Some cheque accounts offer interest on
positive balances. The problem with these bank accounts is that it is
not clear whether the money is being deposited for safe keeping or being
loaned for a fixed term. The distinction between demand deposits and
term deposits has become blurred. In reality they are two different
The two types of account are different for the bank and have
different legal consequences. The bank is entitled to lend money in a
term deposit to someone else for the term of the deposit. It has a duty
to lend the money in away that will minimise risk of loss. The situation
with a demand deposit is very different. If the bank lends money
deposited in a demand deposit, or even records it as an asset in its
accounts, it is stealing something that does not belong to it. To avoid
being guilty of theft, an honest bank should be very clear about what
type of account it is offering.
Understanding the two types of account is also important for the
person making the deposit. When making the demand deposit, they are
giving their money to the bank for safekeeping. Security and convenience
is their priority, so they will be willing to pay a fee for that
service. They also expect to be able to go to the bank and be absolutely
certain that their money will be there.
When a person makes a time deposit, they understand that their money
will no be available to them for the length of the specified term. They
understand that the bank will be lending the money on to someone who can
use it effectively. They expect the bank to be careful and lend the
money wisely. However, they know that there is some risk that the person
borrowing the money may default on the loan. This risk is shared with
other term depositors at the bank and with the owners of the bank, so a
small default will not affect them. However, if defaults become
widespread, the risk to the depositor may increase. The interest paid on
the term deposit is in part compensation for this risk. Generally this
risk premium will be quite small.
Imagine two banks.
Bank No 1 says,
We will store your money and keep it safe. We will not use your money
as if it belonged to us. We will deliver it to any person, as you
instruct us. To cover the cost of providing this service, we will charge
a small monthly fee, but you can be sure that your money will be here
when you want it.
Bank No 2 says,
We will look after your money for you and we will deliver it to
anyone according to your instructions. However, if we see that you are
not using your money, we will lend it to someone who can make use of it.
The interest will compensate you for the risk.
The risk you run is that when you want your money, the person who
borrowed it may not have returned the money back to us. We will have to
borrow the money from some else to repay you, so you may have to wait to
get your money. If for some reason a whole lot of people all decide to
withdraw their money at the same time and there is a run on the bank,
many of you could lose your money. This risk is unlikely and the benefit
is that you will have no fees.
Which bank would people choose?
I do not care which bank people would choose. They can work out the
risks and the benefits and do what is best for them. If they prefer the
cheaper option, they are free to choose it, provide that they do not
expect me to rescue them, if things go wrong.
The problem today is that we do not have a choice. Most people think
that they are dealing with Bank No 1, whereas there bankers believe they
are Bank No 2. There is a dangerous disconnect between the understanding
of risks and expectations.
Before assuming that Bank No 2 is the preferred option, read this
If John returns from overseas and puts $10,000 in his check account,
the balance sheet of the bank shows an increase of $10,000 under cash.
If John were the first client of the bank, its balance sheet would look
It is true that the bank also records a liability to John. However,
because the bank has control of the cash, it has a stronger position.
John is a creditor of the bank, so he is very dependent on the bank
honouring its obligations.
If the bank thinks that John is unlikely to withdraw its money, it
may make a loan to Pete. Its balance sheet would no look like this.
Pete buys a truck from Bill, who deposits the cash he received in the
bank. The bank’s balance sheet now looks like this.
The banks cash to asset ratio is fifty percent, so the bank now
complies with the Basel Accords. It easily meets the standards set by
governments all over the world.
Despite this compliance, we now have a strange situation where the
bank only has $10,000 cash, yet both John and Bill think they have
$10,000 in the bank. If they both try to withdraw their cash out at the
same time, they will not be able to get it. The bank cannot call in the
loan from Pete, because he no longer has the cash. He has brought a
Bill and John’s cash is not lost, but they cannot get hold of it
when they demand it. The best the bank could do is to give $5,000 to
both Bill and John and make them wait for the rest of their money. The
bank could force Pete to repay his loan, but if he has to sell the truck
quickly, he might only get $5,000 for it. John and Bill would then be in
trouble, as one of them would have lost $5,000.
Who Should Decide
Who should decide if money in a cheque account is not being used and
is available for lending to someone else? In the modern system the bank
makes this decision. However, the bank cannot read the minds of its
depositors, so it does not know what they are planning to do with their
money, or when they will want to spend it. Banks can only work on
average behaviour and past experience. Neither are good predictors of
The person who is best placed to decide whether money is available
for lending is the depositor. They know their plans for the money. They
know what they are planning to do. They are best placed to decide how
much they need to hold back to deal with unexpected expense.
If banks only paid interest on term deposits and there was a small
fee for the operation of an account where the money is available on
demand, then people would respond to these incentives. They would
quickly identify money that they do not need immediately and transfer it
into a term deposit that paid interest. The bank would then know that it
was available for lending to others, without having to make guesses
about the depositor’s intentions.
The banking system would function better if banks stopped deciding
when money was available for lending and left these decisions to
Some people think that I should not worry that the bank recording my
deposit as an asset on their balance sheet. They do not under the nature
of modern banking law. When a person deposits money in a bank, they
change from being an owner of an asset into a creditor of the bank. They
give up a property right in exchange for a contractual right. They swap
the ownership of an asset for a promise to receive repayment on demand.
Owning an asset is generally better than being a creditor. If I get a
3 year loan of $30,000 from General Motors, I can buy a new truck. My
balance sheet looks like this.
Before making the sale, the truck was on the General Motors balance
sheet, as part of its inventory. Once they have sold the truck to me, it
becomes my asset. I can drive wherever I like in the truck. I can paint
it bright blue, if I choose. General Motors have lost control of the
truck. They cannot control its use it any more.
General Motors has swapped the truck for a credit contract for the
money that I have borrowed. They cannot demand it back until the three
years are complete. A contractual right is less certain than a property
right. When they owned the truck they knew exactly what they had. The
credit contract is more risky, because they cannot be certain that they
will get their money back when it is due. I may have got into financial
trouble and be unable to repay the loan. The bank has the uncertainty of
a credit contract, whereas I have a property right in the truck. The one
holding the property right is in a better position than the one holding
the credit contract.
In the same way, holding cash provides more security than holding a
credit contract. Being a creditor of the bank is not the same as being
the owner of cash. If I demand repayment and the bank fails to pay me, I
cannot charge it with theft. All I can do is sue the bank for breach on
contract and demand payment of damages. If the bank defaults, my
contract right is converted into a claim in bankruptcy. I have to line
up with other creditors and take my chances.
All that a bank depositor “owns” is the right to enforce the bank
to keep its promise. I believe that most people want better security for
their money. We need banks that offer a different option.
Bailment is an important legal concept.
Bailment is the process of placing personal property or goods in the
temporary custody or control of another. The custodian or holder of the
property, who is responsible for the safe keeping and return of the
property, is know as the bailee. The person who delivers or transfers
the property to the bailee is known as the bailor. For a bailment to be
valid, the bailee must have actual physical control of the property. The
bailee is generally not entitled to the use of the property while it is
in his possession, and a bailor can demand to have the property returned
to him at any time (Lawyers.com)
A bailment is not the same as a sale, which is an intentional
transfer of ownership of personal property in exchange for something of
value. A bailment involves only a transfer of possession or custody, not
of ownership. A bailment is created when a parking garage attendant, the
bailee, is given the keys to a motor vehicle by its owner, the bailor.
The owner, in addition to renting the space, has transferred possession
and control of the vehicle by relinquishing its keys to the attendant (Free
The main distinguishing feature of a bailment is that possession of
the property is transferred to the bailee, but ownership remains with
the bailor. Trusting my furniture to a warehousing company for storage
is an example of a bailment. The warehouse has possession my furniture,
but I am still its owner.
The second feature of a bailment is that the bailee is not entitled
to use the property for their own purpose. The warehouse company is not
entitled to use my furniture.
If a deposit in a cheque account is a bailment, then the bank could
not record the cash as an asset on its balance sheet. It could not use
the cash for its own purposes.
During the 19th century, the British Law Lords ruled that a demand
deposit is not a bailment. This decision has since between adopted by
courts all over the world.
In a case in 1811, Sir William Grant ruled that money paid into a
bank is not a bailment, but a loan. The banker is not a bailee, but a
debtor (Carr v Carr). In a subsequent case, he said, “The money paid
into a banker immediately becomes a part of his general assets and he is
merely a creditor for the amount (Devayne v Noble).
Lord Cottenham summed up the early decisions in
v. Hill and Others.
Money, when paid into a bank, ceases altogether to be the money of
the principal; it is then the money of the banker, who is bound to an
equivalent by paying a similar sum to that deposited with him when he is
asked for it . . . . The money placed in the custody of a banker is, to
all intents and purposes, the money of the banker, to do with it as he
pleases; he is guilty of no breach of trust in employing it; he is not
answerable to the principal if he puts it into jeopardy, if he engages
in a hazardous speculation; he is not bound to keep it or deal with it
as the property of his principal....
According to modern law, a bank deposit is not a bailment, so the
bank is entitled to record the deposit as an asset on its balance sheet.
Bailment not Bailout
That is the situation under law, but there is not reason why the
relationship between a bank and a depositor should be decided by judges.
Banks are providing a service to their customers. The customers are
entitled to demand whatever service they want. If enough customers
demand a particular service, an astute bank would provide that service.
If enough people demanded a different service, some banks would start to
Many people will be happy to accept the service provided by banks
under the current legal arrangements. They are happy for the bank to
take control of their money and use it as they please, provided they can
get good interest and no fees on their cheque account. If that is what
they want, that is fine.
However, many other customers will want a different account. They
will prefer a bailment-type cheque account in which their money does not
get transferred to the balance sheet of the bank, but remains the
property of the depositor. Many customers do not want to become
creditors of the bank. They want to be the owner of their money.
Modern banking law does not prevent banks from contracting with
depositors to provide a bailment-type account. The money in this account
would not become an asset of the bank. It would not be available for the
bank to use. It would be stored in the same way as a storage warehouse
stores my furniture.
When my salary goes into my cheque account, I want it to be available
when I demand it. I do not want to the bank to treat my money as its own
property. I do not want the bank to loan my money out to someone else.
If the bank provides this service well, I am happy to pay a small fee
for that service.
I do not want the bank to decide when I am not going to use my money.
I am capable of doing that myself. If I do not need some of my money for
a time, I will move it into a term deposit, so the bank can lend it out
for a time, but I do not want the bank just deciding it can lend the
money in my cheque account whenever it chooses. I do not want the bank
treating my money in the way described by Lord Cottenham above. I want a
cheque account that is a bailment. If enough customers were to demand
this service, some innovative bank should provide that service.
We should stop talking about bailout of banks and start demanding
bailment from banks.
Clarity and Separation
Banks and their customers should be very clear about the service that
is being offered. The problem is that most modern bank accounts operate
their demand deposit accounts as if they were term deposits. Depositors
have accepted this blurring, because they want interest on their money,
even when it is stored for safe keeping. However, they are allowing
banks to commit theft with their money, in return for an interest
payment. Depositors should be much clearer about what they want from a
bank. Do they want safe keeping or do they want interest?
Banks should keep these two types of transaction quite separate.
Demand deposits should be kept separate from funds that that have been
deposited as term deposits. In practice, it would be better if there
were two types of bank. Some banks would provide a safe-keeping and
payments service. They would do this for a fee. The money should always
be available on demand, so they will not pay interest.
Other banks will become loan brokers. They would pay interest, so
they would only accept term deposits. All those dealing with this type
of bank would understand the risks involved in lending money to other
Distinguishing between true demand deposits term deposits would allow
people to make explicit choices about which service they want. Some will
choose to keep their money safe. Others will choose to put some of their
money in time deposits, so they can earn interest.
Changing the System
Modern banking practices are largely shaped by government policies.
Governments have supported the common practice of moving demand deposits
onto their own balance sheets. We should not count on governments to
resolve this problem. Fortunately, individuals can resolve this problem
without waiting for their government.
We should start asking our banks who owns our demand deposits. They
will tell us that that the money belongs to us. We should respond by
asking them why it appears on their balance sheet, if it belongs to us.
They will respond with “weasel words” about modern banking practice
and banking laws and deposit insurance. However, if enough people
persist in asking these questions, banks will start get concerned.When
there is sufficient demand for a true safekeeping service, some banks
will start to provide it. A bank that has an eye for opportunity will
see an opportunity to get a head start in a niche market.
The customer is king. Once a few banks start offering this service,
people who want a genuine safe-keeping service will switch their
business to those banks. If enough people are serious about honest
banking, the other banks will get worried about losing business. Many
will start offering a true safe-banking service, so that they do not
lose customers. Some will set up separate institutions in attempt a
significant share of this business.
Power to the People
Depositors should also start asking banks what they will do with
their term deposits. They should ask if the money will be loaned for a
longer term than the term of the deposit. If the bank says yes, they
should ask how the bank will repay their money when the term of the
deposit has ended. The bank will say that it will obtain deposits from
other people or borrow the money from other institutions. This kind of
questioning will expose the dangers in what the banks are doing.
If enough people ask for a different service, an innovative bank will
be able to get an advantage by providing that service. If consumers
start enquiring about a bank that matches the terms of it loans with the
terms of its deposits, a bank should emerge to provide that service. If
that is what most depositors really want, then that bank should grow
quickly. As more and more depositors choose this service, other banks
will have to start providing it, so that they do not lose market share.
The power to change the banking system lies with depositors. If
enough people demand a better service, banks will have to change their
practices. All that is need is for one bank to provide an alternative
service. Depositors can then use the power of their money to reward that
bank and punish those that refuse to change. Power rests with those who
own the money and the money is owned by the depositors, not the banks. |
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Space and Time: Inertial Frames
A “frame of reference” is a standard relative to which motion and rest may be measured; any set of points or objects that are at rest relative to one another enables us, in principle, to describe the relative motions of bodies. A frame of reference is therefore a purely kinematical device, for the geometrical description of motion without regard to the masses or forces involved. A dynamical account of motion leads to the idea of an “inertial frame,” or a reference frame relative to which motions have distinguished dynamical properties. For that reason an inertial frame has to be understood as a spatial reference frame together with some means of measuring time, so that uniform motions can be distinguished from accelerated motions. The laws of Newtonian dynamics provide a simple definition: an inertial frame is a reference-frame with a time-scale, relative to which the motion of a body not subject to forces is always rectilinear and uniform, accelerations are always proportional to and in the direction of applied forces, and applied forces are always met with equal and opposite reactions. It follows that, in an inertial frame, the center of mass of a system of bodies is always at rest or in uniform motion. It also follows that any other frame of reference moving uniformly relative to an inertial frame is also an inertial frame. For example, in Newtonian celestial mechanics, taking the “fixed stars” as a frame of reference, we can determine an (approximately) inertial frame whose center is the center of mass of the solar system; relative to this frame, every acceleration of every planet can be accounted for (approximately) as a gravitational interaction with some other planet in accord with Newton's laws of motion.
This appears to be a simple and straightforward concept. By inquiring more narrowly into its origins and meaning, however, we begin to understand why it has been an ongoing subject of philosophical concern. It originated in a profound philosophical consideration of the principles of relativity and invariance in the context of Newtonian mechanics. Further reflections on it, in different theoretical contexts, had extraordinary consequences for 20th-century theories of space and time.
- 1. Relativity and Reference Frames in Classical Mechanics
- 2. Inertial Frames in the 20th Century: Special and General Relativity
- 2.1 Inertial Frames in Newtonian Spacetime
- 2.2 The Conflict Between Galilean Relativity and Modern Electrodynamics
- 2.3 Special Relativity and Lorentz Invariance
- 2.4 Simultaneity and Reference-Frames
- 2.5 From Special Relativity and Lorentz Invariance to General Relativity and General Covariance
- 2.6 The Equivalence of Inertia and Gravity
- 2.7 The Equivalence Principle and General Covariance
- 2.8 The Extension of the Relativity Principle
- 2.9 From Inertial Frames to Curved Spacetime
- Other Internet Resources
- Related Entries
The term “reference frame” was coined in the 19th century, but it has a long prehistory, beginning, perhaps, with the emergence of the Copernican theory. The significant point was not the replacement of the earth by the sun as the center of all motion in the universe, but the recognition of both the earth and the sun as merely possible points of view from which the motions of the celestial bodies may be described. This implied that the basic task of Ptolemaic astronomy—to represent the planetary motions by combinations of circular motions—could take any point to be fixed, and that, as Copernicus suggested in the opening arguments of “On the revolutions of the heavenly spheres,” the choice of any particular point required some justification on other than astronomical grounds. As the basic programme of Ptolemy and Copernicus gave way to that of early classical mechanics, this equivalence of points of view was made more precise and explicit. Galileo demonstrated that the Copernican view does not contradict our experience of a seemingly stable earth, through a principle that, in the precise form that it takes in Newtonian mechanics, has become known as the “principle of Galilean relativity”: mechanical experiments will have the same results in a system in uniform motion that they have in a system at rest. Therefore the experiments claimed as evidence against Copernicus—e.g., that a stone dropped from a tower falls to the base of the tower, instead of being left behind—would happen just as they do whether the earth were moving or not, provided that the motion is sufficiently uniform. See Figure 1.
Figure 1: Galileo's Argument
If the earth is rotating sufficiently uniformly, a stone dropped from the tower will fall straight to the base, just as a stone dropped from the mast of a uniformly moving ship will fall to the foot of the mast. In both cases the stone's vertical motion will be smoothly composed with its horizontal motion. Hence a sufficiently uniform motion will be indistinguishable from rest.
Leibniz, later, articulated a more general “equipollence of hypotheses”: in any system of interacting bodies, any hypothesis that any particular body is at rest is equivalent to any other. Therefore neither Copernicus' nor Ptolemy's view can be true—though one may be judged simpler than the other—because both are merely possible hypothetical interpretations of the same relative motions. This principle clearly defines (what we would call) a set of reference frames, differing in their arbitrary choices of a resting point or origin, but agreeing on the relative positions of bodies at any moment and their changing relative distances through time.
For Leibniz and many others, this general equivalence was a matter of philosophical principle, founded in the metaphysical conviction that space itself is nothing more than an abstraction from the geometrical relations among bodies. In some form or other it was a widely shared tenet of the 17th-century “mechanical philosophy”. Yet it was flatly incompatible with physics as Leibniz himself, and the other “mechanists,” actually conceived it. For the basic program of mechanical explanation depended essentially on the concept of a privileged state of motion, as expressed in the assumption that bodies maintain a state of rectilinear motion until acted upon by an external cause. Thus their fundamental conception of force, as the power of a body to change the state of another, likewise depended on this notion of a privileged state. This dependence was clearly exhibited in the vortex theory of planetary motion, in which every orbit was explained by the balance between the planet's inherent centrifugal tendency (its tendency to follow the tangent to the orbit) and the pressure of the surrounding medium.
For this reason, the notion of a dispute between “relativists” or “relationists” and “absolutists” or “substantivalists”, in the 17th century, is a drastic oversimplification. Newton, in his controversial Scholium on space, time, and motion, was not merely asserting that motion is absolute in the face of the mechanists' relativist view; he was arguing that a conception of absolute motion was already implicit in the views of his opponents—that it was implicit in their conception, which he largely shared, of physical cause and effect. The general equivalence of reference-frames was implicitly denied by a physics that understood forces as powers to change the states of motion of bodies.
Newton therefore held that physics required the conception of absolute space, a distinguished frame of reference relative to which bodies could be said to be truly moving or truly at rest. Assuming, as both Newton and Leibniz did, that states of motion could be distinguished by their causes and effects, the distinguished status of this frame of reference is physically well founded—and metaphysically well-founded for a metaphysics that, like Newton's or Leibniz's, takes force to be a well-founded notion. On Leibniz's conception of force, in particular, a given force is required to generate or to maintain a given velocity—for objects “passively” resist motion, but maintain their states of motion only by “active” force—so that, on dynamical grounds, “every body truly does have a certain amount of motion, or, if you will, force.” This implies that there is in principle a distinguished frame of reference in which the velocities of bodies correspond to their true velocities, i.e. to the amounts of moving force that they truly possess, and it implies that in any frame that is in motion relative to this one, bodies will not have their true velocities. In short, such a conception of force, if it could be applied physically, would give a precise physical application of Newton's conception of absolute space.
The difficulty with Newton's view of absolute space comes from the Newtonian conception of force. If force is defined and measured solely by the power to accelerate a body, then obviously the effects of forces—in short, the causal interactions within a system of bodies—will be independent of the velocity of the system in which they are measured. So the existence of a set of equivalent “inertial frames” is imposed from the start by Newton's laws. Suppose that we determine for the bodies in a given frame of reference—say, the rest frame of the fixed stars—that all observable accelerations are proportional to forces impressed by bodies within the system, by equal and opposite actions and reactions among those bodies. Then we know that these physical interactions will be the same in any frame of reference that is in uniform rectilinear motion relative to the first one. Therefore no Newtonian experiment will be able to determine the velocity of a body, or system of bodies, relative to absolute space. In other words, there is no way to distinguish absolute space itself from any frame of reference that is in uniform motion relative to it. Newton thought that a coherent account of force and motion requires a background space consisting of “places” that “all keep given positions in relation to one another from infinity to infinity” (1726, p. 412). But the laws of motion enable us to determine an infinity of such spaces, all in uniform rectilinear motion relative to each other, and furnish no way of singling out any one as “immovable space.”
Oddly enough, no one in the 17th century, or even before the late 19th century, expressed this equivalence of reference-frames more clearly than Newton himself. Newton explicitly derived it from the laws of motion as Corollary V:
When bodies are enclosed in a given space, their motions in relation to one another are the same whether the space is at rest or whether it is moving uniformly straight forward without circular motion. (1726, p. 423.)
This is the first clear statement of the Galilean relativity principle. It implied that the dispute between the heliocentric and geocentric views of the universe was mistakenly framed: the proper question about “the system of the world” was not “which body is at rest in the center?” but “where is the center of gravity of the system, and which body is closest to it?” For in a system of orbiting bodies, only their common center of gravity will be unaccelerated, and by Corollary V, the motions of the bodies in the system will be the same, whether its center of gravity is at rest or in uniform rectilinear motion. The system is indeed approximately Keplerian, since the sun has by far the greatest mass and is therefore little disturbed from the center of gravity, which is therefore very close to the common focus of the approximately Keplerian ellipses in which the planets orbit the sun. But by Corollary V, the nearly-Keplerian structure of the system is completely independent of the system's state of motion in absolute space.
The Galilean relativity principle thus expressed the insight that different states of uniform motion, or different uniformly-moving frames of reference, determine only different points of view on the same physically objective quantities, namely force, mass, and acceleration. We can see this insight expressed more explicitly in Newton's understanding of inertia. For Leibniz (among others) , as we saw, moving force, the power of a body to change the motion of another, was determined by velocity. It was therefore seen as an active power, fundamentally different from the passive power of a resting body to resist any change of position. Newton, in contrast, understood the “force of inertia” as a Galilei-invariant quantity:
[A] body exerts this force only during a change of its state, caused by another force impressed upon it, and the exercise of this force is, depending on viewpoint, both resistance and impetus: resistance in so far as the body, in order to maintain its state, strives against the impressed force, and impetus in so far as the same body, yielding only with difficulty to the force of a resisting obstacle, endeavors to change the state of that obstacle. Resistance is commonly attributed to resting bodies and impetus to moving bodies; but motion and rest, in the popular sense of the term, are distinguished from each other only by point of view, and bodies commonly regarded as being at rest are not always truly at rest. (1726, p. 404–05.)
Newton thus recognized the powers distinguished by Leibniz as the same thing seen from different points of view.
Newton understood the Galilean principle of relativity with a degree of depth and clarity that eluded most of his “relativist” contemporaries. It may seem bizarre, therefore, that the notion of inertial frame did not emerge until more than a century and a half after his death. He had identified a distinguished class of dynamically equivalent “relative spaces,” in any of which true forces and masses, accelerations and rotations, would have the same objectively measured values. Yet these spaces, though empirically indistinguishable, were not equivalent in principle; evidently Newton conceived them as moving with various velocities in absolute space, though those velocities could not be known. Why should not he, or someone, have recognized the equivalence of these spaces immediately?
This is not the place for an adequate answer to this question, if indeed one is possible. For much of the 20th century, the accepted answer was that of Ernst Mach: Newton lived in an age “deficient in epistemological critique,” and so was unable to draw the conclusion that these empirically indistinguishable spaces must be equivalent in every meaningful sense, so that no one of them deserves even in principle to be designated as “absolute space.” Yet even those whom the 20th century credited with more sophisticated epistemological views, such as Leibniz, evidently had difficulties understanding force and inertia in a Galilei-invariant way, despite a philosophical commitment to relativity. Perhaps it suffices to say that to abandon the intuitive association of force or motion with velocity in space, and to accept an equivalence-class structure as the fundamental spatiotemporal framework, requires a level of abstraction that became possible only with the extraordinary development of mathematics, especially of a more abstract view of geometry, that took place in the 19th century. (See geometry: in the 19th century.) In the 17th century only Christiaan Huygens came close to expressing such a view; he held that not velocity, but velocity-difference, was the fundamental dynamical quantity. He therefore understood, for example, that the “absoluteness” of rotation had nothing to do with velocity relative to absolute space, but arose from the difference of velocity among different parts of a rotating body—a difference which would, evidently, be the same irrespective of the velocity of the body as a whole in absolute space. But of this Huygens gave only the merest suggestion, in manuscripts that remained unpublished for two centuries. (See Stein 1977.) The concept of inertial frame therefore emerged only in the late 19th century, when, as we shall see, it did not seem to be of any great immediate importance.
Meanwhile, the relativity principle was understood as the equivalence of uniform states of motion, but any systems in such a state was implicitly understood to have a definite, though unknown and unknowable, velocity in absolute space. Euler (1748), for example, defended Newton's conceptions of space and time against the thesis that space and time are ideal, and motion merely relative; his broad argument was that metaphysics had no standing to criticize conceptions that are required by the established laws of physics. Yet he noted that the laws of motion permit us to determine, not the velocity of any motion in space, but only the absolute sameness of direction of an inertial trajectory over time, and the equality of time-intervals in which an inertially-moving particle moves equal distances. To Euler, these irreducibly spatial and temporal aspects of the laws of motion implied that space and time could not possibly be ideal. Like Newton, therefore, he upheld both the relativity of velocity and the reality of absolute space. The inconsistency of such a theory can be seen in two ways. On the one hand, we can see it as a fundamental incoherence, even if, again, we excuse those who held it on the grounds of the limited mathematical tools available to them. On the other hand, it does represent a deep appreciation of the indistinguishability of velocities in absolute space, and a consequent effort to make sure that the actual treatment of actual physical systems is not undermined by this uncertainty. Newton hoped to analyze the dynamical interactions that hold the solar system together; he wanted to show that his dynamical account, and the view of “the frame of the system of the world” that emerges from it, is a matter of “reasoning from phenomena” rather than of plausible conjecture. It was therefore a very circumspect, even prescient, move on his part to demonstrate, through his use of Corollaries IV and V, that the analysis is completely independent of any conceivable translation of the system in absolute space.
The development of this concept began with a renewed critical analysis of the notion of absolute space, for reasons not anticipated by Newton's contemporary critics. Its starting point was a critical questions about the law of inertia: relative to what is the motion of a free particle uniform and rectilinear? If the answer is “absolute space,” then the law would appear to be something other than an empirical claim, for no one can observe the trajectory of a particle relative to absolute space. Two quite different answers to the question were offered in 1870, in the form of revised statements of the law of inertia. Carl Neumann proposed that when we state the law, we must suppose that there is a body somewhere in the universe—the “body Alpha”—with respect to which the motion of a free particle is rectilinear, and that there is a time-scale somewhere relative to which it is uniform (Neumann 1870). Ernst Mach (1883) claimed that the law of inertia, and Newton's laws generally, implicitly appeal to the fixed stars as a spatial reference-frame, and to the rotation of the earth as a time-scale; at least, he held, such is the basis for any genuine empirical content that the laws have. The notion of absolute space, it followed, was only an unwarranted abstraction from the practice of measuring motions relative to the fixed stars.
Mach's proposal had the advantage of a clear empirical motivation; Neumann's “body Alpha” seemed no less mysterious than absolute space, and almost sounds comical to the modern reader. But Neumann's discussion of a time-scale was somewhat more fruitful. He noted that the law of inertia defines a time-scale: equal intervals of time are those in which a free particle travels equal distances. Such a definition is another aspect of the Newtonian theory first made explicit by Euler (1748). Neumann also noted, however, that this definition is quite arbitrary. For, in the absence of a prior definition of equal times, any motion whatever can be stipulated to be uniform. It is no help to appeal to the requirement of freedom from external forces, since the free particles presumably are known to us only by their uniform motion. We have a genuine empirical claim only when we state of at least two free particles that their motions are mutually proportional; equal intervals of time can then be defined as those in which two free particles travel mutually proportional distances.
Neumann's definition of a time-scale directly inspired Ludwig Lange's conception of “inertial system,” introduced in 1885 . An inertial coordinate system ought to be one in which free particles move in straight lines. But any trajectory may be stipulated to be rectilinear, and a coordinate system can always be constructed in which it is rectilinear. And so, as in the case of the time-scale, we cannot adequately define an inertial system by the motion of one particle. Indeed, for any two particles moving anyhow, a coordinate system may be found in which both their trajectories are rectilinear. So far the claim that either particle, or some third particle, is moving in a straight line may be said to be a matter of convention. We must define an inertial system as one in which at least three non-collinear free particles move in noncoplanar straight lines; then we can state the law of inertia as the claim that, relative to an inertial system so defined, the motion of any fourth particle, or arbitrarily many particles, will be rectilinear. The notions of inertial system and Neumann's time-scale, which Lange called an “inertial time-scale,” may be combined as follows: relative to a coordinate system in which three free particles move in straight lines and travel mutually-proportional distances, the motion of any fourth free particle will be rectilinear and uniform. The questionable Newtonian concepts of absolute rotation and acceleration, Lange proposed, could now be replaced by the concepts of “inertial rotation” and “inertial acceleration,” i.e. rotation and acceleration relative to an inertial system and inertial time-scale. See Figures 2 and 3.
Figure 2: Neumann's Time-Scale:
By Newton's first law, a particle not subject to forces travels equal distances in equal times. But which particles are free of forces? This might appear to be a matter of convention.
Either P1 or P2 can be arbitrarily stipulated to be at the origin of a system of coordinates, and to serve as the measure of equal times But I can say of two particles with different velocities: in intervals of time in which one moves a given distance d1, the other moves a proportional distance d2 = kd1 (where k is a constant; i.e., d1/d2 = k). Or I can compare a particle to a freely rotating planet: in intervals of time through which the planet rotates through equal angles, the particle moves equal distances.
Figure 3: Lange's Definition of ‘inertial system’ (1885):
An inertial system is a coordinate system with respect to which three free particles, projected from a single point and moving in non-coplanar directions, move in straight lines and travel mutually-proportional distances. The law of inertia then states that relative to any inertial system, any fourth free particle will move uniformly.
At about the same time, apparently unaware of the work of Mach, Neumann, and Lange, James Thomson expressed the content of the law of inertia, and the appropriate frame of reference and time-scale (“dial-traveller”), somewhat differently:
For any set of bodies acted on each by any force, a REFERENCE FRAME and a REFERENCE DIAL-TRAVELLER are kinematically possible, such that relatively to them conjointly, the motion of the mass-centre of each body, undergoes change simultaneously with any infinitely short element of the dial-traveller progress, or with any element during which the force on the body does not alter in direction nor in magnitude, which change is proportional to the intensity of the force acting on that body, and to the simultaneous progress of the dial-traveller, and is made in the direction of the force. (Thomson 1884, p. 387)
More simply, an inertial reference-frame is one in which Newton's second law is satisfied, so that every acceleration corresponds to an impressed force. Thomson did not reject the term “absolute rotation,” holding instead that it has to be understood as rotation relative to a reference frame that satisfies his definition. The definition does not express, as Lange's does, the degree of arbitrariness involved in the construction of an inertial system by means of free particles. Moreover, like Lange's, it leaves out a crucial condition for an inertial system as we understand it: all forces must belong to action-reaction pairs. Otherwise we could have, as on a rotating sphere, merely apparent (centrifugal) forces that are, by definition, proportional to mass and acceleration, and so the rotating sphere would satisfy Thomson's definition. Therefore the definition needs to be completed by the stipulation that to every action there is an equal and opposite reaction. (This completion was actually proposed by R.F. Muirhead in 1887.)
But, so completed, Thomson's definition has two advantages over Lange's. First, by appealing to Newton's second law instead of his first, it shows that we can apply the notion of inertial frame without having to consider the question whether there really are any free particles in nature. Second, it exhibits more clearly an essential point about the relation between the laws of motion and the inertial frames: that the laws assert the existence of at least one inertial frame. The original question, “relative to what frame of reference do the laws of motion hold?” is revealed to be wrongly posed. For the laws of motion essentially determine a class of reference frames, and (in principle) a procedure for constructing them. For the same reason, a skeptical question that is still commonly asked about the laws of motion—why is it that the laws are true only relative to a certain choice of reference frame?—is also wrongly posed. If Newton's laws are true, then we can construct an inertial frame; their truth doesn't depend on our ability to construct such a frame in advance.
By the early years of the 20th century, this notion of inertial system seems to have been widely accepted, even if the specific works of Lange and Thomson were largely forgotten; in writing “On the electrodynamics of moving bodies” in 1905, Einstein took it to be obvious to his readers that classical mechanics does not require a single privileged frame of reference, but an equivalence-class of frames, all in uniform motion relative to each other, and any of which “the equations of mechanics hold good.” Two inertial frames with coordinates (x, y, z, t) and (x′, y′ z′ t′) are related by the Galilean transformations,
x′ = x − vt
y′ = y
z′ = z
t′ = t
(assuming that the x axis is defined to be the direction of their relative motion). These transformations clearly preserve the invariant quantities of Newtonian mechanics, i.e. acceleration, force, and mass (and therefore time, length, and simultaneity). As far as Newtonian mechanics was concerned, then, the problem of absolute motion was completely solved; all that remained was to express the equivalence of inertial frames in a simpler geometrical structure.
The lack of a privileged spatial frame, combined with the obvious existence of privileged states of motion—paths defined as rectilinear in space and uniform with respect to time—suggests that the geometrical situation ought to be regarded from a four-dimensional spatiotemporal point of view. The structure defined by the class of inertial frames can be captured in the statement that spacetime is a four-dimensional affine space, whose straight lines (geodesics) are the trajectories of particles in uniform rectilinear motion. See Figure 4.
Figure 4: Inertial Trajectories as Straight Lines of Spacetime
The uniformly moving particle will travel the same distance in the same intervals. A particle that accelerates after t1 will move a greater distance during t2 and therefore its path in spacetime changes direction.
That is, spacetime is a structure whose automorphisms—the Galilean transformations that relate one inertial frame to another—are equivalent to affine transformations: they take straight lines into straight lines (i.e. an inertial motion in one inertial frame will be an inertial motion in any other inertial frame, and likewise for an accelerating or rotational motion), and parallel lines into parallel lines (i.e. uniformly-moving particles or observers who are relatively at rest in one frame will also be relatively at rest in another). (See Stein 1967, Ehlers 1973, and Friedman 1983 for further explanation.) An inertial frame can be characterized as a family of parallel straight lines “filling” spacetime, representing the possible trajectories of a family of free particles that are relatively at rest. See Figure 5:
Each of these families of straight lines, F1 and F2, represents the trajectories of a family of free particles that are relatively at rest, and therefore each defines an inertial frame. Relative to each other, the frames defined by F1 and F2 are in uniform motion.
Each of the surfaces S is a “hypersurface of absolute simultaneity” representing all of space at a given moment; evidently (given the Galilean transformations) two inertial frames will agree on which events in spacetime are simultaneous.
From this we can see that the assertion that an inertial frame exists imposes a global structure on spacetime; it is equivalent to the assertion that spacetime is flat. As we can see from the Galilean transformations, distinct inertial frames will agree on time and simultaneity. Therefore, in the four-dimensional picture, the decomposition of spacetime into hypersurfaces of absolute simultaneity is independent of the choice of inertial frame. Another way of putting this is that Newtonian spacetime is endowed with a projection of spacetime onto time, i.e. a function that identifies spacetime points that have the same time-coordinate. Similarly, absolute space arises from a projection of spacetime onto space, i.e. a function that identifies spacetime points that have the same spatial coordinates. See Figure 6.
|The relation of simultaneity “decomposes” spacetime into 3-dimensional pieces, each representing “all of space at a given time,” by projecting spacetime onto time, i.e., by identifying spacetime points that have the same time coordinates.||Similarly, one can think of the notion of “same place” as projecting spacetime onto space, i.e., by identifying spacetime points that have the same spatial coordinates; each of the trajectories thus singled out represents “a given place at all times.”|
But this latter projection is arbitrary: while it assumes that we can identify the same time at different spatial locations, Newtonian mechanics provides no physical way of identifying the same spatial point at different times. Thus the equivalence of inertial frames can be thought of as the arbitrariness of the projection of spacetime onto space, any such projection being, essentially, the arbitrary choice of some particular inertial frame as a rest-frame.
|Here is a spacetime diagram of motions relative to the inertial frame in which O1, O2, and P are at rest. This can be seen as arising from the projection of each of their inertial trajectories onto a single point of space.||Here is the same situation viewed from an inertial frame in which O3 and P′ are at rest. Now O1, O2, and P are in uniform motion.|
|O1 and O2 are at rest||O3 is at rest|
|O3 is in uniform motion||O1 and O2 are in uniform motion|
|O4 is accelerating any old way||O4 is accelerating any old way|
|O5 and O6 are revolving around their common centre of gravity P, which is at rest||O5 and O6 are revolving around their common centre of gravity P, which is in uniform motion|
|O7 and O8 are revolving around their centre of gravity P′, which is in uniform motion.||O7 and O8 are revolving around their centre of gravity P′, which is at rest|
By the time that this representation of the Newtonian spacetime structure was developed, however, the Newtonian conception of inertial frame had been essentially overthrown. First, 19th-century electrodynamics raised again the question of a privileged frame of reference: the conception of light as an electromagnetic wave in the ether implied that the rest-frame of the ether itself should play a distinguished role in electrodynamical phenomena. On the one hand, physicists such as Maxwell and Lorentz were careful to point out that velocity relative to the ether was not equivalent to absolute velocity, and that the state of motion of the ether itself was necessarily unknown—in other words, that this conception of light did not violate the classical principle of relativity. On the other hand, the existence of such a preferred frame made the equivalence of inertial frames correspondingly less interesting, even if it was true in principle. This is why the appearance of the idea of inertial frame in the 1880's, as I suggested earlier, was not of pressing physical interest to the majority of physicists, and seemed to be a mere philosophical sidelight. The attempts to measure the effects of motion relative to the ether commanded considerably more attention.
Second, the abandonment of the ether—following the failure of attempts to measure velocity relative to the ether and, more generally, the apparent independence of all electrodynamical phenomena of motion relative to the ether—did not vindicate the Newtonian inertial frame, but required a dramatically revised conception. Special relativity might be said to have applied the relativity principle of Newtonian mechanics to Maxwell's electrodynamics, by eliminating the privileged status of the rest-frame of the ether and admitting that the velocity of light is independent of the motion of the source. As Einstein expressed it, “the same laws of electrodynamics and optics will be valid for all frames of reference for which the equations of mechanics hold good.” (1905, p. 38.) But as Einstein also pointed out, the invariance of the velocity of light and the principle of relativity, at least in its Galilean form, are incompatible. It simply makes no sense, according to Galilean relativity, that any velocity should appear to be the same in inertial frames that are in relative motion.
Einstein solved this difficulty through his analysis of simultaneity: frames in relative motion can agree on the velocity of light only if they disagree on simultaneity; only the relativity of simultaneity makes possible the invariance of the velocity of light. This means that the transformations between inertial frames that preserve the velocity of light will not preserve simultaneity. These are the Lorentz transformations:
Evidently these transformations do not preserve length and time, and so the invariant quantities of Newtonian mechanics, which presuppose invariant measures of length and time, must now depend on the choice of inertial frame. By the same token, the notions of force, mass, and acceleration can no longer be appealed to in the definition of an inertial frame. The definition must instead appeal to the invariant quantities of electrodynamics: an inertial frame is one in which light travels equal distances in equal times in arbitrary directions. What seems impossible, from the point of view of Galilean relativity, is that a frame that moves uniformly relative to such a frame should also satisfy the definition. But that, again, rests on the assumption that two inertial frames will have a common measure of simultaneity. If, as Einstein asserts, the only reasonable definition of simultaneity is one provided by light signals, then there is no determination of simultaneity that will give the same results in different inertial frames. The spacetime structure that is implied by special relativity is thus an affine space, like Newtonian spacetime, but it is not objectively divided into hypersurfaces of absolute simultaneity; the sets of simultaneous events for any inertial frame are the hyperplanes orthogonal to the trajectories that determine that frame. In other words, the choice between two inertial frames determines a choice between two distinct divisions of spacetime into space and time. See Figure 8:
The inertial frames F and F′ are in relative motion, and therefore, as the Lorentz transformations indicate, they disagree on simultaneity. F and F′ thus determine distinct decompositions of spacetime into instantaneous spaces, S and S′, respectively
The details of Einstein's argument and the structure of Minkowski spacetime can be found elsewhere (see, e.g., Einstein 1951 and Geroch 1978). Here only one more point is worth making. It could be argued that Einstein's and Lorentz's view are completely equivalent. That is, we could assume that there is indeed a privileged frame of reference, and that the apparent invariance of the velocity of light is explained by the effects on bodies of their motion through the ether (the Lorentz contraction and time dilation). This purported distinction between empirically indistinguishable frames has often been criticized on straightforward methodological grounds, but it could be (and surely has been) argued that it is more intuitively plausible than the relativity of simultaneity. After all, knowing that (as Einstein showed) the Lorentz contraction can be derived from the invariance of the velocity of light does not, by itself, entitle us to say which of the two is the more convincing starting-point.
This is why it is so important that Einstein's 1905 paper begins with a critical analysis of the entire notion of a frame of reference. It is tacitly assumed by Lorentz's theory, and classical electrodynamics generally, that we have a reference-frame in which we can measure the velocity of light. But how is such a reference-frame determined? The distances between points in space can only be determined if it is possible to determine which events are simultaneous. In practice this is always done by light-signalling, if only in the informal sense that we identify simultaneous events when we see them at the same time. But if the spatial frame of reference is determined by light-signals, and is then to be used to measure the speed of light, we would appear to be going in a circle; the underlying assumption must be that, while light-signalling is useful and practical, it is not essential to the definition of simultaneity, and that there is a fact of the matter about which events are simultaneous that is independent of this method of signalling. This assumption was actually made explicit by James Thomson. He recognized—alone, apparently, before Einstein—that the measurement of distance involves
the difficulty as to imperfection of our means of ascertaining or specifying, or clearly idealizing, simultaneity at distant places. For this we do commonly use signals by sound, by light, by electricity, by connecting wires or bars, and by various other means. The time required in the transmission of the signal involves an imperfection in human powers of ascertaining simultaneity of occurrences at distant places. It seems, however, probably not to involve any difficulty of idealizing or imagining the existence of simultaneity. Probably it may not be felt to involve any difficulty comparable to that of attempting to form a distinct notion of identity of place at successive times in unmarked space. (1884, p. 380).
In other words, Thomson assumed that it was not a difficulty in principle, like the difficulty of determining rest in absolute space. But Einstein showed that it was precisely the same kind of difficulty, and that determinations of simultaneity involve reference to an arbitrary choice of reference-frame, just as much as determinations of velocity. Einstein's conclusion is, of course, entirely contingent on the empirical facts of electrodynamics; it could have been avoided if there were in nature a useful signal of some kind whose transmission would provide a criterion of absolute simultaneity, so that the same events would be determined to be simultaneous in all inertial frames. Or, experiments might have been able to reveal the dependence of the velocity of light on the state of motion of the source. Then synchronization by light-signals could still have been regarded as a mere practical substitute for a notion of absolute simultaneity that stood on independent grounds, empirically as well as conceptually. But as Einstein saw, because of the apparent independence of the velocity of light of the motion of the source, even “idealizing or imagining the existence of simultaneity” involves light-signaling more essentially than anyone could have realized. Unless some other criterion of simultaneity is provided, therefore, the establishment of a spatial frame of reference involves light-signaling in an essential way. In the absence of such a criterion the speed of light cannot be, as Lorentz supposed, empirically measured against the background of an inertial frame; in that case the only empirically sound definition of an inertial frame is the one that appeals to the speed of light.
It may seem surprising that, after this insightful analysis of the concept of inertial frame and its role in electrodynamics, Einstein should have turned almost immediately to call that concept into question. But he had a compelling combination of physical and philosophical motives to do so. On the physical side, he realized (along with many others) that special relativity would require some fundamental revision of the Newtonian theory of gravity. On the philosophical side, he became convinced, largely by his reading of Mach (1883), that the central role of inertial frames was an “epistemological defect” that special relativity shared with Newtonian mechanics. (Einstein 1916, pp. 112–113.) Only relative motions are observable, yet both of these theories purport to identify a privileged state of motion and use it to explain observable effects (such as centrifugal forces). Coordinate systems are not observable, yet both of these theories assign a fundamental physical role to certain kinds of coordinate system, namely, the inertial systems. In either theory, inertial coordinates are distinguished from all others, and the laws of physics are said to hold only relative to inertial coordinate systems. In an epistemologically sophisticated theory, both of these problems would be solved at once: the new theory would only refer to what is observable, which is relative motion; it would admit arbitrary coordinate systems, instead of confining itself to a special class of system. Why, after all, should any genuine physical phenomenon depend on the choice of coordinate system?
Another way of putting the same point is to say that, in Newtonian mechanics and special relativity, rotation is “absolute” because the transformations between inertial frames (Galilean or Lorentzian) preserve rotational states. Thus the “absoluteness” of rotation arises precisely from singling out one type of frame, by one type of transformation, instead of allowing arbitrary transformations and arbitrary frames. Einstein held that this epistemological insight had a natural mathematical representation in the principle of general covariance, or the principle that the laws of nature are to be invariant under arbitrary coordinate transformations. More precisely, what this means is that coordinate transformations are no longer required (as in the affine spaces of Newtonian mechanics and special relativity) to take straight lines to straight lines, but only to preserve the smoothness of curves (i.e. their differentiability). The general theory of relativity was intended to be a generally covariant account of spacetime, and its general covariance was intended to express the general relativity of motion. And the theory came into being because Einstein perceived a deep connection between this project and that of finding a relativistic theory of gravitation.
The philosophical motivations and implications of Einstein's view are dealt with elsewhere. (See, for example, the entries on Einstein's philosophy of science; the hole argument; and early philosophical interpretations of general relativity.) We will consider here only the bearing of general relativity on the notion of an inertial frame. It is questionable whether Einstein succeeded in establishing the general relativity of motion, but it is clear that general relativity undermines the concept of inertial frame in important respects. This arises from the equivalence principle: that inertial mass—the quantity that enters into Newton's second law, and that is a measure of a body's resistance to acceleration—is equivalent to gravitational mass, the quantity that enters into Newton's law of universal gravitation. A more empirical way of expressing it is that all bodies fall with the same acceleration in the same gravitational field, or, the trajectory of a body in a given gravitational field will be independent of its mass and composition. This is the principle that Newton tested by constructing pendulums with wooden boxes as their bobs, which he would fill with different materials in order to see whether those differences made a difference to the speed of falling; they didn't. Eötvös made more precise tests in the late 19th century, and established the principle to much greater accuracy; these are the results on which Einstein would have relied. Newton also tested the principle for bodies whose masses differ greatly, by observing that Jupiter and its four moons all received precisely the same acceleration from the sun's gravitational field.
The equivalence principle suggests, however, that a freely-falling frame of reference is physically indistinguishable from an inertial frame. Newton had already noticed this, and indeed he stated it, more or less, in Corollary VI to the laws of motion:
If bodies are moving in any way whatsoever with respect to one another and are urged by equal accelerative forces along parallel lines, they will all continue to move with respect to one another in the same way as they would if they were not acted on by those forces. (1726, p. 423.)
For example, he was able to treat the system of Jupiter and its moons as if it were (nearly) at rest or moving uniformly in a straight line, because the attractive force of the sun acts (almost) equally on every part of the system. See Figure 9:
Figure 9: Newton's Corollary VI
What seem, within a given system, like equal and parallel accelerations may be, on a larger scale, unequal and converging on some distant massive object; e.g., the system of Jupiter and its moons is falling toward the sun, but “locally” the accelerations are very nearly equal and parallel, and may therefore be neglected.
He even applied this reasoning to the entire solar system, in order to justify treating it as an isolated system: if there were any outside force acting on it, it must have been acting more or less equally and in parallel directions on all parts of the system.
It may be alleged that the sun and planets are impelled by some other force equally and in the direction of parallel lines; but by such a force (by Cor. VI of the Laws of Motion) no change would happen in the situation of the planets to one another, nor any sensible effect follow; but our business is with the causes of sensible effects. Let us, therefore, neglect every such force as imaginary and precarious, and of no use in the phenomena of the heavens….(1729, volume 2 p. 558)
Now, it is a familiar fact that in an orbiting spacecraft, bodies behave as if no forces were acting on any of them (as if they were “weightless”), because the attraction of the earth acts equally on all of them. But these phenomena are not, by themselves, evidence that no phenomena are capable of distinguishing an inertial frame from a falling frame. Einstein was willing to generalize the equivalence principle, and to conclude that the classical idea of a distinguished class of frames of reference has no physical basis. Any frame that we might regard as inertial might be, for all we can tell by experiment, in free fall. By the same token, any frame that is uniformly accelerating is indistinguishable from one that is at rest in a uniform gravitational field. Suppose that you are in a box at rest on the earth; you and everything in the box, by the equivalence principle, will be accelerated downward with the acceleration g (= 9.8 meters/second/second). Now suppose that the box itself is in empty, gravity-free space, but accelerating upward (i.e. in the direction of its roof) with the acceleration -g. Obviously, because of their inertia, bodies in the box, including your own, will exert the same force—have the same “weight”—on the floor as if the box were at rest and sitting on the earth.
To get a clearer idea of the physical significance of the equivalence principle, and its connection with general covariance, consider the Newtonian procedure for analyzing motion in the solar system, here sketched very roughly:
- Determine the accelerations of all the planets relative to the fixed stars.
- Using the laws of motion, their corollaries, and all the propositions proved from these in Book I of Principia, derive from the accelerations the forces needed to produce them; in particular, derive from the orbits the centers of those orbits, and the masses of the bodies needed to produce those forces. This crucially involves the law of action and reaction, for otherwise it would be impossible to break down the total acceleration of any planet into the components contributed by particular other planets; the earth's acceleration, for example, is the sum of its accelerations toward all the other planets, and each individual acceleration is part of an action-reaction pair involving some other planet.
- When we understand the mutual interactions among the planets, we are in a position to estimate their relative masses. In Newton's case, this was necessarily restricted to the planets with satellites, because only in those cases could he compare the accelerations they determine at given distances and so deduce the differences in mass. By this reasoning he estimated the ratios of the Sun's mass to those of Jupiter (1067 to 1), Saturn (3021 to 1), and the earth, and was able to calculate that the center of mass of the entire solar system would never be more than one solar diameter from the center of the Sun.
- Having found the center of mass, we have in principle determined an inertial frame: by Corollary IV to the laws of motion, the center of mass will be at rest or moving uniformly in a straight line. That is, the mutual actions of the bodies in the system will not change the state of motion of the center of mass. And having determined an inertial frame, we are in a position to say that the accelerations relative to the center of mass frame are the true accelerations.
One might think that the problem of relativity arises right from the start: the reliance on the fixed stars already seems to introduce an arbitrary assumption that threatens to vitiate Newton's procedure as an account of the true motions. But the framework of the fixed stars, initially just taken for granted, turns out to be justified in the course of the analysis. If it turns out that all the accelerations relative to the fixed stars can be analyzed into action-reaction pairs involving bodies within the system, leaving no “leftover” accelerations that need to be traced to some yet-unknown influence, then we can conclude that the stars are a suitable (sufficiently inertial) frame of reference after all. (By the later 19th century, observations became sufficiently precise to reveal that there is in fact a leftover acceleration, namely the famous extra precession of Mercury. But that could not affect Newton's analysis in 1687.) In contrast, had we chosen the earth as a frame of reference, we would find that there are accelerations relative to this frame—e.g. Coriolis and centrifugal accelerations—that don't satisfy the law of action and reaction.
The relativistic aspect of this situation arises from the equivalence principle. Newton's Corollary VI said that the inertial frame we construct by this procedure is effectively indistinguishable from one in which all the bodies are undergoing equal and parallel accelerations caused by some force that acts equally on all of them; the equivalence principle asserts that gravity is such a force. In following the Newtonian procedure for constructing an inertial frame, we have constructed a frame which might be, for all we can determine empirically, falling in the gravitational field of some other system. Here again, as in his use of Corollary V, we can see that Newton was being remarkably circumspect about his frame of reference: he needed to show that his analysis of the forces at work, and his conclusion about the nearly-heliocentric structure of the system, are not affected by any unknown forces acting on the system as a whole, and his appeal to Corollary VI precisely satisfies this need. By the same token, however, the accelerations relative to this frame cannot be known to be the “true accelerations”; they may be accelerations relative to a freely-falling trajectory just in case the center of mass is itself freely falling, in which case they have to be added to the gravitational acceleration of the center of mass before we can arrive at the true accelerations. But the acceleration of the center of mass may have to be added to some larger acceleration—and so on. This means that we can't know the true strength of the gravitational field by observing the motions in this frame. The only hope of doing so would be to include all the mass in the universe in one dynamical system; if we knew the center of mass of the entire universe, we could rule out the possibility that something else is exerting an accelerative force, since by hypothesis there would be nothing else.
We can see the significance of this more clearly by looking at the equations of motion (in a very simplified form). Newton's equation of motion for a particle subject to no force asserts that it moves uniformly, with zero acceleration. Obviously, in a gravitational field, the particle's acceleration will depend on the field. In effect, we are accounting for the trajectory of the falling particle by “decomposing” it into two parts, the part determined by its natural tendency to move uniformly in a straight line, and the part contributed by the gravitational field. But by the analysis of the equivalence principle, determining the inertial part—and therefore determining the gravitational part—depends on our assumption that the center of mass frame is inertial rather than freely falling. And this assumption is arbitrary; that is, it amounts to an arbitrary choice of the coordinate system in which we define the equation of inertial motion. This implies that the gravitational field depends on the coordinate system in precisely the same way.
The principle of general covariance, then, acquires its physical significance in conjunction with the equivalence principle. By itself, it says that the geometrical structures of spacetime don't depend on the coordinates in which we express them, or on the set of points that we may think comprises spacetime. This is an important principle, but it doesn't recommend general relativity over other theories, since special relativity and Newtonian mechanics also involve spacetime structures that can be defined in a generally-covariant way, through the same kinds of coordinate-independent mathematical objects that we use in general relativity. Combined with the equivalence principle, however, it implies that a central Newtonian idea—that gravity is a force causing deviations from uniform rectilinear motion—is based on an arbitrary choice of coordinates. For a trajectory that satisfies all empirical criteria for being inertial in a particular frame of reference—e.g. the trajectory of the center of mass in our example—may be freely falling relative to some other trajectory that satisfies the same criteria. By contrast, a freely-falling trajectory is a freely falling trajectory in any coordinate system; it is only the decomposition of it into its inertial and gravitational parts that will be different in different coordinate systems.
General covariance is thus not an argument against privileged states of motion, as Einstein had hoped it would be. It is an argument that the privileged states of motion should not be mere artifacts of our choice of coordinates, i.e. that they should be coordinate-independent. Precisely what this means depends, then, on what physical means we have at our disposal to identify states of motion other than by simply setting down coordinates. Combined with the equivalence principle, it is an argument for regarding gravitational free-fall as the privileged state of motion, rather than as a forced deviation from the privileged state of motion. And in this way it provides an argument for spacetime curvature. As we saw, in Newtonian and Minkowski spacetime the inertial trajectories are, by definition, the straight lines or geodesics of spacetime. And the flatness of spacetime consists in the fact that these geodesics behave like straight lines in a flat space or surface: parallel geodesics remain parallel, and non-parallel geodesics do not accelerate relative to one another. (In any inertial frame, the motion of any other inertial frame appears uniform.) By the equivalence principle, however, free-fall trajectories satisfy all empirical criteria for being inertial trajectories, and so the distinction between the two types of trajectory depends on the mere choice of coordinates. General covariance suggests, then, that the free-fall trajectories ought to be identified as the inertial trajectories—and therefore, as the geodesics of spacetime. But if free-fall trajectories are the geodesics of spacetime, then spacetime is curved. For the free-fall trajectories exhibit relative accelerations, and the relative acceleration of geodesics is a defining characteristic of curved geometry. The curvature of the earth's surface, for example, is revealed in the fact that geodesics that begin in parallel directions can begin to approach one another—for example, two lines of longitude can both be perpendicular to the equator, but converge on one another as they approach the poles. And since the relative accelerations of falling bodies depend on the distribution of mass, as we already knew from Newton's theory, we now conclude not only that spacetime is curved, but that its curvature is determined by the distribution of mass. (For further explanation see Geroch 1978.)
The curvature of spacetime, finally, determines the status of inertial frames in general relativity. The statement that all reference-frames, rather than just inertial frames, are equivalent is a misleading way of describing the situation; rather, the variable curvature of spacetime makes the imposition of a global inertial frame impossible. So the status of the latter is like the status of a plane rectangular coordinate system on the surface of the earth. Over a sufficiently small area, the coordinate plane may be a good approximation to the surface, but over increasingly large areas it diverges increasingly from the contours of the earth. And if two such coordinate systems, with their origins at different points on the earth, are extended until they meet, they will be seen to be “disoriented” relative to one another. In contrast, a flat plane can be so coordinatized, and coordinate systems originating at different points can be smoothly combined into one system. Similarly, in the affine spaces of Newtonian and special-relativistic physics, any inertial coordinate system can be extended over the whole of spacetime. And in any system so extended, the trajectory of every other inertial observer will be a uniform rectilinear motion. But if spacetime is variably curved, according to the distribution of mass and energy, local inertial systems will be “disoriented” relative to one another; indeed, the degree of this “disorientation” is one of the measures of curvature. And an inertially-moving—i.e. freely falling—particle will in that case be accelerating in the local inertial system of another freely-falling particle. Thus there are inertial trajectories, but no extended inertial systems. See Figures 10–13:
This Cartesian coordinate system can evidently be simply “set down” over the plane below. Any coordinate system defined at any point of the plane can be smoothly extended over the entire plane.
Figure 11: “Magnified” View of Flat “Local” Coordinate Systems on a Curved Surface
This arbitrary curved surface won't allow for the global laying down of a coordinate system, but must be coordinatized in small overlapping pieces, which generally won't be parallel to one another.
In a flat spacetime, the rest-frame of any inertial observer an be “extended” over all of spacetime in such a way that, in this global inertial frame, the trajectory of every other inertial observer will be an inertial trajectory.
In a curved spacetime, inertial trajectories will be relatively accelerated; indeed the relative acceleration of geodesics is a measure of curvature. Therefore the local inertial frame of any freely-falling observer cannot be extended into a global frame in which all other inertial observers are moving uniformly. The inertial frames of different freely-falling observers will be, like local coordinate systems on a curved surface, “disoriented” relative to one another.
One could try to express this idea with Einstein's remark about the need to “free oneself from the idea that coordinates must have an immediate metrical meaning.” (Einstein 1949, p. 67.)But even this might be misleading. Einstein evidently was thinking that, in general relativity, coordinates, and coordinate transformations, no longer represent the possible displacements of rigid bodies or the transport of ideal clocks. The insight underlying this is that the notion of rigid displacement—therefore of rigid coordinate system, and inertial frame—imposes a priori a degree of uniformity, or symmetry, on spacetime; the displacement of bodies without change of dimension, and the transport of an ideal clock without distortion of time-intervals, requires a homogeneous space. And so rigid displacement cannot be a basic principle in a theory in which spacetime curvature varies according to the distribution of mass and energy. The possibility of a rigid displacement, and therefore the existence of an inertial frame, can only arise a posteriori, as the result of a peculiar distribution of mass-energy (for example, in a universe empty of mass and energy, or with a highly symmetrical distribution). The serious defect in the notion of inertial frame is not that it makes an arbitrary distinction among coordinate systems—for the distinction is quite as genuine as the distinction between flat and curved spacetime—but that it extends indefinitely over spacetime a structure that, in our universe, only corresponds approximately to very small regions.
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Einstein, Albert: philosophy of science | general relativity: early philosophical interpretations of | geometry: in the 19th century | space and time: conventionality of simultaneity | space and time: the hole argument |
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