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-{"text": "The following includes information that may be useful in understanding the present invention(s). It is not an admission that any of the information provided herein is prior art, or material, to the presently described or claimed inventions, or that any publication or document that is specifically or implicitly referenced is prior art."}
-{"text": "1. Field of the Invention\nThis invention relates to frequency and award redemption program. More particularly, the present invention relates to an on-line, interactive frequency and award redemption program which is fully integrated.\n2. Description of Related Art\nFrequency programs have been developed by the travel industry to promote customer loyalty. An example of such a program is a \u201cfrequent flyer\u201d program. According to such a program, when a traveler books a flight, a certain amount of \u201cmileage points\u201d are calculated by a formula using the distance of the destination as a parameter. However, the mileage points are not awarded until the traveler actually takes the flight.\nWhen a traveler has accumulated sufficient number of mileage points, he may redeem these points for an award chosen from a specific list of awards specified by the program. Thus, for example, the traveler may redeem the points for a free flight ticket or a free rental car. In order to redeem the accumulated points, the traveler generally needs to request a certificate, and use the issued certificate as payment for the free travel.\nWhile the above program may induce customer loyalty, it has the disadvantage that the selection of prizes can be made only from the limited list of awards provided by the company. For example, a traveler may redeem the certificate for flights between only those destinations to which the carrier has a regular service. Another disadvantage is that the customer generally needs to plan ahead in sufficient time to order and receive the award certificate.\nAccording to another type of frequency and award program, a credit instrument is provided and credit points are accumulated instead of the mileage points. In such programs, bonus points are awarded by using a formula in which a price paid for merchandise is a parameter. Thus, upon each purchase a certain number of bonus points are awarded, which translate to dollar credit amount. According to these programs, the customer receives a credit instrument which may be acceptable by many enrolled retailers, so that the selection of prizes available is enhanced. An example of such a program is disclosed in E.P.A. 308,224. However, while such programs may enhance the selection of prizes, there is still the problem of obtaining the credit instrument for redeeming the awarded points. In addition, the enrollee must allow for processing time before the bonus points are recorded and made available as redeemable credit. Thus, the immediacy effect of the reward is lacking in these conventional incentive programs."}
-{"text": "In axial piston machines, at least one working piston is mounted in a longitudinally displaceable manner in a cylinder bore of a piston drum and forms a cylinder space with the cylinder bore. The cylinder space is alternately compressed and decompressed by the longitudinal movement of the working piston and, accordingly, alternately connected to a high-pressure reservoir and a low-pressure reservoir. When changing from the low-pressure reservoir connection to the high-pressure reservoir connection, pulsations occur which can result in substantial noise generation. To counteract this, so-called pre-compression volumes are used, which are formed by pre-compression spaces.\nAxial piston machines having pre-compression spaces or zones are known from, the prior art, for example from DE 197 06 114 C5. In this, a pre-compression volume or a reservoir element is integrated in a control plate or in a connection plate of the axial piston machine. The pre-compression volumes known from the prior art can additionally be controlled via valve devices.\nThe known pre-compression spaces are sealed or closed outside the housing of the axial piston machine, which requires additional installation space. With regard to automobile applications, the installation space is an increasingly important issue for axial piston machines.\nThe seal of the pre-compression spaces moreover poses a technical challenge owing to the pulsation.\nThe object of the disclosure is to reduce the installation space for creating pre-compression spaces."}
-{"text": "The present invention relates generally to optical communications and, more particularly, to multiple symbol polarization switching differential-detection modulation formats.\nAs Internet traffic grows exponentially because of a variety of user terminals and internet services, it has prompted strong research interests on high-speed optical networks, which are the backbone infrastructure of current \u201cGlobe Village\u201d. The data rate for optical fiber communications has moved from 10 Gbits/s to 40 Gb/s and 100 Gbits/s or even 1 Tbits/s per channel. However, one of the major challenges facing the ultra-high-data-rate dense wavelength division multiplexing (DWDM) optical fiber transmissions is the fiber nonlinearity, which causes optical signal distortions due to the various nonlinear effects in optical fiber and sets the limit of the maximal reach. DQPSK modulation is an important format for high-speed optical communications by transmitting 2 bits per symbol. At 40 Gb/s, DQPSK systems employing direct detection are attractive by having low complexity and being generally available.\nIn a digital coherent optical communication system, different types of digital signal processing (DSP) functions can be applied, to mitigate the fiber nonlinearity, such as digital back-propagation algorithms. However, the existing DSP-based fiber nonlinearity mitigation algorithms are demanding on the hardware resources, which are relatively limited and sophisticated due to the requirements of very high electronic processing bandwidth. Meanwhile, most of the existing nonlinearity mitigation algorithms show very limited system performance improvements in real experimental testing.\nIn another approach, the phase conjugation scheme has been proposed to improve the systems' nonlinearity tolerance. However, the deployment of this scheme requires at the exact middle point of the entire transmission link, thus imposing a strict and thus unpractical restrictions on the system deployment. Its spectrum efficiency would be halved because of the fiber four-wave mixing effects. In another prior effort, the polarization states for adjacent symbols are arranged in orthogonal states for improving fiber nonlinearity tolerance.\nAccordingly, there is a need for a low-cost solution to increase the nonlinearity tolerance of a direct-detection optical DQPSK system"}
-{"text": "A. Field of the Invention\nThe present invention relates to digging implements, and in particular a backhoe.\nB. Background of the Art\nBackhoes are used extensively for excavating and for carrying objects from one area to another. Backhoes have typically been used to dig holes in the ground for trenches and for the placement of building structural components, road substructures, cables, pipes, etc.\nHeretofore, backhoes have included a first arm pivotally attached to a tractor and a second arm pivotally attached to the first arm in a scissors-like manner. A bucket is attached to the second arm for digging. Separate hydraulic actuators have typically been used to move each of the arms and the bucket. Some of these backhoes have included an extendable second arm. Furthermore, some backhoes have included a gripping device positioned opposite the bucket for gripping objects between the gripping device and the bucket. One of the gripping devices has included a gripping device statically attached to an arm of the backhoe that does not rotate relative to the arm. These gripping devices have been difficult to use because the arm and the bucket have to properly position relative to the gripping device before the gripping device can be used to pick up objects. Another gripping device includes a separate hydraulic actuator for moving only the gripping device. These backhoe are expensive to manufacture because of the cost for the extra hydraulic actuator and the cost for connecting the gripping device to the controls in the tractor. A third gripping device includes thumbs that rotate simultaneously with the bucket. These backhoes are also difficult to use because the bucket and the gripping device must be properly positioned before the gripping implement can be used. Furthermore, these backhoes are difficult to operate because the rotating gripping implement can get in the way of the rotating bucket, thus making the ground difficult to dig."}
-{"text": "The phytopathogenic fungus Ashbya gossypii is a filamentously growing ascomycete that was first isolated as a plant pathogen in tropical and sub-tropical regions. It infects the seed capsule of cotton plants (Ashby S. F. and Nowell W. (1926) Ann. Botany 40: 69-84) and has also been isolated from tomatoes and citrus fruits (Phaff H. J. and Starmer W. T. (1987) In \"The Yeasts\", Vol. I Rose A. H., Harrison, J. S. (eds), Academic Press, London, 123 ff; Dammer K. H. and Ravelo H. G. (1990). Arch. Phytopathol. Pflanzenschutz, Berlin 26: 71-78 Dammer and Ravelo, 1990). The infection of the seed capsule is caused by transmission of A. gossyppii mycelium pieces or spores by stinging-sucking insects and causes a disease called stigmatomycosis.\nStudies characterising the karyotype of A. gossypii have been performed (Wright, 1990; Wendland, 1993; Gaudenz, 1994, \"The small genome of the filamentous fungus Ashbya gossypii: Assessment of the karyotype\", Diploma Thesis, Department of Applied Microbiology, Biocenter, University Basel). It has been found using yeast chromosomes of precisely known length as size markers that the genome of A. gossypii has a total nuclear genome size of 8.85 Mb. Presently, A. gossypii represents the most compact eukaryotic genome, compared to genome sizes of 12.5 Mb for Saccharomyces cerevisiae (Chu et al. (1986) Science, 234:1582-1585), 31.0 Mb for Aspergillus nidulans (Brody and Carbon (1989) Proc Natl Acad Sci USA. 86:6260-6263), and 47.0 Mb for Neurospora crassa (Orbach et al.(1988) Mol Cell Biology, 8:1469-1473).\nA. gossypii is systematically grouped to the endomycetales belonging to the family of spermophthoraceae. This classification is based on the observation that the spores that develop in hyphal compartments called sporangia look like ascospores, which are defined as end products of meiosis.\nSince A. gossypii is a filamentous ascomycete, and is capable of growing only by filamentous (hyphal) growth, fungal targets found in this model organism are predictive of targets which will be found in other pathogens, the vast majority of which grow in a filamentous fashion."}
-{"text": "Yo-yo players, especially beginners, have been assisted by the development of yo-yos provided with a means to automatically return the yo-yo to the player's hand before the yo-yo spins out completely. Such an arrangement is described, for example, in U.S. Pat. No. 4,332,102 to Caffrey.\nThe Caffrey patent discloses a yo-yo having a rotatable bearing pulley mounted on the axle to which the yo-yo string is attached. Adjacent the pulley section of the bearing there is provided a cylindrical friction or braking means that interacts with two clutch mechanisms. The surface of the friction or braking means has a slip resistant characteristic and is in practice one or a series of O-rings, which are subject to wear. The clutch mechanism is provided with weighting means such that when the yo-yo is thrown the clutch is released by the development of centrifugal forces. The centrifugal forces are counterbalanced by a spring-loaded force such that the clutch is activated when the yo-yo slows down. The clutch engages the cylindrical friction surface of the pulley extension while the yo-yo still has sufficient momentum to enable the automatic return of the yo-yo to the player's hand. The successful development of an automatically returning yo-yo has proven to be especially valuable to beginners. It is also viewed as a valuable assistance to less gifted or handicapped players.\nThe nature of the arrangement shown in Caffrey, however, is such that it tends to critically weaken the structural integrity of the yo-yo. The pulley bearing to gain access to the clutch mechanism housed within a yo-yo half requires part of the boss enclosing the axle to be removed in the plastic mold. To enable sufficient braking capacity to be applied, up to 80% of the plastic boss must be removed where the pulley extension friction surface meets the clutch mechanisms.\nAlso, because the pulley is also the string bearing means, problems occur when bearing lubrication applied in excess finds its way with the aid of centrifugal forces to the nearby contact area between clutch and pulley such that the clutch slips and fails to return the yo-yo successfully.\nThe arrangement disclosed in Caffrey, by linking the centrifugally operated clutch to a coaxial extension of the string bearing means, has intrinsically restricted options on varying the quality of the string bearing means. The use of a dual purpose bearing that combines a string securing means as well as a clutch interfacing means where the clutch means is operatively enclosed in the yo-yo half must by nature expand laterally along the axial member to accommodate both functions. Caffrey achieves a superior spinning automatically returning yo-yo is achieved by narrowing the string bearing means thereby reducing the area frictionally contacting the axial member. Having the centrifugally activated clutch operatively engage an integral extension of the string bearing means also necessitates the use of double-loop string attachment to inhibit the string from slipping on the bearing means and thereby reducing the clutch effectiveness. The general public have difficulty in tying a double-loop attachment.\nAlso, the Caffrey arrangement, at least in its commercial embodiments in which the clutch mechanism occupies only one yo-yo half, exhibits a weight differential between the two yo-yo halves that is believed to shorten free spinning time.\nAnother problem experienced with a conventional automatically returning yo-yo which has a static spacing between yo-yo halves is that different yo-yo manoeuvres, to be performed efficiently, require different yo-yo response tolerances. Tom Kuhn in his publication \"SB2Flight Manual\" in \"The Art of Yo-Yo Choreography\", indicates a narrower string gap is better for loop-the-loops and a wider string gap is better for complex spin tricks."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to the field of electrically-driven reciprocating pumps. More particularly, the invention relates to a pump which is particularly well suited for use as a fuel pump, driven by a solenoid assembly employing a permanent magnet and a solenoid coil to produce pressure variations in a pump section and thereby to draw into and express from the pump section a fluid, such as a fuel being pumped. The invention also relates to a fuel injector assembly employing such a pump.\n2. Description of the Related Art\nA wide range of pumps have been developed for displacing fluids under pressure produced by electrical drives. For example, in certain fuel injection systems, fuel is displaced via a reciprocating pump assembly which is driven by electric current supplied from a source, typically a vehicle electrical system. In one fuel pump design of this type, a reluctance gap coil is positioned in a solenoid housing, and an armature is mounted movably within the housing and secured to a guide tube. The solenoid coil may be energized to force displacement of the armature toward the reluctance gap in a magnetic circuit defined around the solenoid coil. The guide tube moves with the armature, entering and withdrawing from a pump section. By reciprocal movement of the guide tube into and out of the pump section, fluid is drawn into the pump section and expressed from the pump section during operation.\nIn pumps of the type described above, the armature and guide tube are typically returned to their original position under the influence of one or more biasing springs. Where a fuel injection nozzle is connected to the pump, an additional biasing spring may be used to return the injection nozzle to its original position. Upon interruption of energizing current to the coil, the combination of biasing springs then forces the entire movable assembly to its original position. The cycle time of the resulting device is the sum of the time required for the pressurization stroke during energization of the solenoid coil, and the time required for returning the armature and guide to the original position for the next pressure stroke.\nWhere such pumps are employed in demanding applications, such as for supplying fuel to combustion chambers of an internal combustion engine, cycle times can be extremely rapid. Moreover, repeatability and precision in beginning and ending of pump stroke cycles can be important in optimizing the performance of the engine under varying operating conditions. While the cycle time may be reduced by providing stronger springs for returning the reciprocating assembly to the initial position, such springs have the adverse effect of opposing forces exerted on the reciprocating assembly by energization of the solenoid. Such forces must therefore be overcome by correspondingly increased forces created during energization of the solenoid. At some point, however, increased current levels required for such forces become undesirable due to the limits of the electrical components, and additional heating produced by electrical losses.\nThere is a need, therefore, for an improved technique for pumping fluids in a linearly reciprocating fluid pump. There is a particular need for an improved technique for providing rapid cycle times in fluid pumps, such as fuel pumps without substantially increasing the forces and current demands of electrical driving components."}
-{"text": "1. Field of the Invention\nThe present invention is directed to a low pressure plasma generator with a localizable plasma combustion chamber.\n2. Discussion of the Background\nTreatment with low pressure plasmas is an important new method for modifying the surfaces of solid bodies. The surfaces can be, e.g., etched, i.e., partially removed, or activated, i.e., in an energy-rich state that is suitable for extensive modifications, or are coated by bonding gaseous substances. For all of these methods, the surface to be modified must be subjected to a plasma. As is well-known, a gas comprising excited molecules, radicals or ions is referred to as a plasma.\nPlasmas can be generated at low gas pressures by means of microwave radiation. A prerequisite for the formation of a plasma is an adequately high field strength of the radiation. However, the field strength is the greatest in the immediate vicinity of the source of radiation and decreases with increasing distance therefrom. Therefore, the plasma may exist only in the vicinity of the source of radiation.\nThe uniform treatment of large surfaces or surfaces with complicated shapes with a plasma causes considerable difficulties. For reasons relating to their design and their energy supply, available sources of radiation cannot be disposed at any point and at any position in a low pressure chamber. Similarly, the surface to be treated cannot be moved to specified locations in the plasma combustion chamber. Therefore, the surfaces to be treated cannot be located near the plasma source.\nThe ability to ignite and maintain a plasma at a predetermined place, where it is supposed to unfold its technological effect, is called localization. The precise localization of plasma is of great importance primarily when a large surface is to be treated uniformly. This goal can be largely reached if one can succeed in localizing a plasma linearly and moving the plasma uniformly over the surface to be treated. For this purpose, either the plasma can be localized stationarily and the substrate can be moved relative thereto or the substrate can remain stationary and the plasma is moved at right angles to its longitudinal extension. However, just the linear localization of a uniform plasma causes considerable difficulties.\nThe literature reports on various possibilities for localizing microwave plasmas. These include, among others, the ignition of the plasma behind the inlet window for the microwave (Wertheimer et al., Thin Solid Films, 115 (1984), 109), the ignition of primary transmitting aerials (Alcatel DVM, 92240 Malakoff, France, machine GIR 820), the ignition by means of local pressure differences in a vacuum chamber (IKV reports, Mr. Ludwig) and the magnetic confinement with or without the utilization of an electron cyclotron resonance absorption (EP-A 279 895). Some of these possibilities were also used for localizing large area plasmas.\nThe use of surface waveguide structures, which are mounted outside the vacuum apparatus but which are in front of a microwave permeable window, allows a large area plasma to be ignited (Kieser et al., Thin Solid Films, 118 (1984), 203).\nAll of the described methods of localization have drawbacks that stand in the way of their practical application. The drawback of the arrangement described last is that the plasma burns only directly behind the window and cannot be moved within the vacuum to any arbitrary place therein by the operator. In the case of a coating plasma, the window is also coated, a feature that can lead to an absorption and reflection of the microwaves depending on the properties of the deposited layer. Long setting-up times then become necessary owing to the repeated cleaning or exchanging of the windows.\nAn ignition at a primary transmitting aerial yields a plasma whose intensity in most cases exhibits local inhomogeneities owing to the wavelength of the transmitting frequency (e.g., with a period of 12 cm at a frequency of 2.45 GHz). In this arrangement, compensating devices such as a mechanical movement of the aerial can hardly be used owing to the design of primary transmitting aerials\nThe ability to localize a plasma by means of local pressure differences is limited to the coating of largely closed bodies. This is a suitable method for coating bottles internally. However, in trying to process a flat substrate with such a system grave technological problems arise.\nOne successful method is magnetic confinement. This process is used, e.g, in the sputter technique. However, an effective magnetic confinement in achieved only if the gyration radius of the charged particles in the plasma with respect to the free path cannot be ignored. This is the case for conventional permanent magnets made of ferrite only below pressures of about 0.1 mbar.\nWith plasmas of higher pressures--of up to a few millibars--higher etching and deposition rates can be obtained during the etching and coating process. For this reason there is a need for a method for plasma confinement that also has a good localizability at higher pressures and thus allows homogeneous etching or formation of layers at simultaneously high etching and deposition rates."}
-{"text": "1. Field of the Invention\nThe present invention relates to an electronic calculator which is capable of operating a numerical expression in sequence of touching the keyboard thereof in accordance with the order from left to right reading along the expression to be calculated, and is capable, in case of an expression including parentheses, of visually indicating in the display unit thereof a temporary answer resultant from operating the portion of the equation between the parentheses. The present invention also relates to an electronic calculator, which visually indicates not only the temporary result derived from a part of an expression between parentheses but any temporary result obtained from an independently operable portion of an expression to be calculated, by discriminating as an arithmetic block any independently operable portion of the expression to execute in turn an arithmetic operation to that portion.\n2. Description of the Prior Art\nA conventional desktop calculator has been designed with giving importance to the simple system configuration thereof, so that some operational functions might be reduced to a certain extent. For example, when a conventional calculator is operated by touching keys 3, .times., ( , 4, +, 5, ), and = in accordance with a numerical expression 3 .times. (4 + 5) =, the calculator will display or print out only the final result 27, and not any intermediate temporary result derived from a portion of the expression, such as 9 obtained from (4 + 5).\nHowever it is often necessary for an operator to be informed of a temporary result with regard to a portion of an expression to be calculated. Otherwise, an operator must redundantly operate the calculator to obtain the temporary result. In the example described above, when an operator intends to know the result of (4 + 5) as a partial result of the expression, the operator should depress keys 4, +, 5, and = to be informed of the answer 9 visually indicated in the display unit thereof, which answer may be written down on a sheet of paper, and after that the operator will clear all previous settings in the calculator to carry out the remaining operation 3 .times. 9 = by touching keys 3, .times., 9, and =. Then the final answer 27 will be indicated in the display unit. Thus, in a conventional calculator, an operator must discretely twice operate the keyboard thereof in accordance with two numerical expressions such as 4 + 5 = and 3 .times. 9 =. In case of calculating a number of numerical expressions, an operator will be worried by obtaining a lot of the temporary results.\nCalculating a numerical expression containing an exponential term, such as 3.sup.2, an operator often intends to known the partial answer, such as 9 in the aforesaid example, resultant from the exponential operation. In that case, the first touching of keys 3, a.sup.x, 2, and = causes the answer 9 to be obtained in a conventional calculator, and then the remaining part of that expression being calculated by employing the intermediate result 9. Thus, a numerical expression containing many independently operable terms should be divided into portions to be partially calculated so as to obtain intermediate result.\nSince such a conventional calculator is capable of displaying only the final answer from a numerical expression, it becomes more difficult to check out misoperations in keying before completing the calculation of the expression including more terms. Such a conventional calculator requires relatively more careful operation in keying, causing an operator to be exhausted."}
-{"text": "An electromagnetic fuel injector comprises a cylindrical tubular body displaying a central feeding channel, which functions as a fuel conduit and ends with an injection nozzle regulated by an injection valve controlled by an electromagnetic actuator. The injection valve is provided with a needle, which is rigidly connected to a mobile keeper of the electromagnetic actuator in order to be displaced by the action of the electromagnetic actuator between a closed position and an open position of the injection nozzle against the bias of a spring which tends to hold the needle in the closed position. The valve seat is defined in a sealing element, which is shaped as a disc, lowerly and fluid-tightly closes the central channel of the support body and is crossed by the injection nozzle.\nPatent application EP1635055A1 describes an electromagnetic fuel injector in which a guiding element rises from the sealing element, such guiding element having a tubular shape, accommodating the needle therein in order to define a lower guide of the needle itself and displaying a smaller external diameter with respect to the internal diameter of the feeding channel of the supporting body so as to define an external annular channel through which pressurised fuel flows. Four through feeding holes, which lead towards the valve seat to allow the flow of pressurised fuel towards the valve seat itself, are obtained in the lower part of the guiding element. The needle ends with an essentially spherical shutter head, which is adapted to fluid-tightly rest against the valve seat and slidingly rests on an internal cylindrical surface of the guiding element so as to be guided in its movement. The injection nozzle is of the \u201cmulti-hole\u201d type, i.e. it is defined by a plurality of through injection holes, which are obtained from a chamber formed downstream of the valve seat; in this way, the optimal geometries of the injection nozzle may be obtained for the various applications by appropriately orienting the single injection holes.\nExperimental tests have shown that the drive time-injected fuel quantity curve (i.e. the law linking the drive time to the quantity of injected fuel) of the electromagnetic injector described above is on the whole rather linear, but displays an initial step (i.e. displays a step increase for short drive times and therefore for small quantities of injected fuel); furthermore, the extent of such initial step is higher proportionally to the fuel feeding pressure.\nConsequently, the electromechanical injector described above may be used in a direct injection internal combustion Otto cycle engine (i.e. fed with petrol, LPG, methane or the like), in which the fuel feeding pressure is limited (lower than 200-250 bars) and the injector is not normally driven to inject small amounts of fuel). However, the electromagnetic injector described above cannot be used in a small direct injection internal combustion Diesel cycle engine (i.e. fed with Diesel fuel or the like), in which the feeding pressure of the fuel is rather high (up to 800-900 bars) and the injector is constantly driven so as to perform a series of pilot injectors before a main injection."}
-{"text": "The present invention generally relates to a device which eliminates wind rushing noise between a crash helmet and face shield structures.\nWith the ever increasing popularity of relatively high speed motorcycles, conventional protective helmets, while satisfactory to a certain degree, do not satisfy the requirements of all users of such equipment. One type of helmet commonly used is a type which protects the face of the motorcycle operator by providing a partial cylindrical transparent face shield. The periphery of the face shield overlaps the helmet shell and is fixedly secured thereto by a plurality of fasteners such as rivets or any suitable type of releasable fastener such as a screw threaded fastener, snap fastener or the like which would enable the face shield to be removed or replaced in the event of damage thereto. The bottom edge of the face shield and the crash helmet surrounding the rider's upper neck are open so as to provide access to the interior of the helmet and face shield assembly to facilitate it being placed on the head of the wearer and removed therefrom.\nWith the most common types of crash helmets equipped with transparent face shields available today, a gap or space up to one-half (1/2) inch in width exists between the periphery of the face shield where it overlaps and is affixed to the helmet with fasteners, and the underlying front leading edge of the helmet itself. As the cyclist proceeds forward, wind rushing through this gap passes down over the face and ears as it exits out of the bottom of the helmet and affixed face shield surrounding the upper neck. This wind flow disturbance behind the face shield and surrounding helmet results in eye irritation and annoying noise which increases in intensity as the cyclist goes faster and faster. The present invention eliminates this eye irritating and noisy wind flow disturbance while still allowing a milder airflow to promote comfort and help prevent fogging of the face shield."}
-{"text": "1. Field of the Invention\nThe present invention relates to a level gauge for detecting a level of liquid helium which is accommodated in a container made of metals, glasses or other materials. More particularly, the invention relates to a level gauge for detecting a level of liquid helium which uses, as a sensing element, a wire made of an amorphous superconductive alloy obtained by rapid by quenching a molten alloy. The level of liquid helium is detected by measuring an electric current, voltage and/or electric resistance of the sensor element.\n2. Description of Related Technology\nLiquid helium level gauges that use superconducting alloy sensing elements rely on the electrical resistance changes of the element to indicate liquid level. The portion of the element submerged in the liquid becomes superconductive, i.e. no resistance to electrical current. The portion above the liquid is not superconductive and resists electrical flow at a constant rate over its length. If the sensing element is homogeneous, has a constant width, and has a constant thickness, the resistance properties will be constant over the length of the sensor element. By passing an electrical current through the submerged element, measuring the electrical current, and comparing the value to a calibration relationship, the level of the helium can be determined.\nIn the level gauge disclosed in U.S. Pat. No. 4,655,079 (which is herein incorporated by reference), the superconductive alloy is represented by the following formula: EQU Zr.sub.100-x (Ru.sub.y Rh.sub.1-y)x\nwherein x represents the contents of Ru and/or Rh in aomic % and has a numeral value of 22.5<x<27.5; and y represents a numerical value of 0<y<1.\nHowever, the superconducting transition temperature (Tc) of that superconductive alloy ranges from 4.2K to 4.5K. This transition temperature is quite close to the temperature of liquid helium (4.2K). As the pressure in the storage vessel changes, the accuracy of the level measurements can decrease."}
-{"text": "Recently the use of enzymes, especially of microbial origin, has become more and more common. Enzymes are used in several industries including, for example, the starch industry, the dairy industry, and the detergent industry. It is well known in the detergent industry that the use of enzymes, particularly proteolytic enzymes, has created industrial hygiene concerns for detergent factory workers, particularly due to the health risks associated with dustiness of the available enzymes.\nSince the introduction of enzymes into the detergent business, many developments in the granulation and coating of enzymes have been offered by the industry. See for example the following patents relating to enzyme granulation:\nU.S. Pat. No. 4,106,991 describes an improved formation of enzyme granules by including within the composition undergoing granulation, finely divided cellulose fibers in an amount of 2-40% w/w based on the dry weight of the whole composition. In addition, this patent describes that waxy substances can be used to coat the particles of the granulate.\nU.S. Pat. No. 4,689,297 describes enzyme containing particles which comprise a particulate, water dispersible core which is 150-2,000 microns in its longest dimension, a uniform layer of enzyme around the core particle which amounts to 10%-35% by weight of the weight of the core particle, and a layer of macro-molecular, film-forming, water soluble or dispersible coating agent uniformly surrounding the enzyme layer wherein the combination of enzyme and coating agent is from 25-55% of the weight of the core particle. The core material described in this patent includes clay, a sugar crystal enclosed in layers of corn starch which is coated with a layer of dextrin, agglomerated potato starch, particulate salt, agglomerated trisodium citrate, pan crystallized NaCl flakes, bentonite granules or prills, granules containing bentonite, Kaolin and diatomaceous earth or sodium citrate crystals. The film forming material may be a fatty acid ester, an alkoxylated alcohol, a polyvinyl alcohol or an ethoxylated alkylphenol.\nU.S. Pat. No. 4,740,469 describes an enzyme granular composition consisting essentially of from 1-35% by weight of an enzyme and from 0.5-30% by weight of a synthetic fibrous material having an average length of from 100-500 micron and a fineness in the range of from 0.05-0.7 denier, with the balance being an extender or filler. The granular composition may further comprise a molten waxy material, such as polyethylene glycol, and optionally a colorant such as titanium dioxide.\nU.S. Pat. No. 5,254,283 describes a particulate material which has been coated with a continuous layer of a non-water soluble, warp size polymer. U.S. Pat. No. 5,324,649 describes enzyme-containing granules having a core, an enzyme layer and an outer coating layer. The enzyme layer and, optionally, the core and outer coating layer contain a vinyl polymer.\nWO 91/09941 describes an enzyme containing preparation whereby at least 50% of the enzymatic activity is present in the preparation as enzyme crystals. The preparation can be either a slurry or a granulate.\nWO 97/12958 discloses a microgranular enzyme composition. The granules are made by fluid-bed agglomeration which results in granules with numerous carrier or seed particles coated with enzyme and bound together by a binder.\nHowever, even in light of these developments offered by the industry (as described above) there is a continuing need for low-dust granules. In particular, it is especially problematic in the detergent industry when granules in general, or those comprising proteins or enzymes, form dust and are aerosolized. In these cases, workers are often exposed to the contents of the granules and can develop severe allergic reactions. Therefore, it is an object of the present invention to provide a method of producing a low-dust enzyme granule by adding antifoam agent. It is a further object of the invention to facilitate a safer environment for workers in the detergent industry who are exposed to enzyme containing granules."}
-{"text": "Solid gas sorption systems are used to produce cooling and/or heating by repeatedly desorbing and absorbing the gas on a coordinative complex compound formed by absorbing a polar gas refrigerant on a metal salt in a sorption reaction sometimes referred to as chemisorption. Complex compounds incorporating ammonia as the polar gaseous refrigerant are especially advantageous because of their capacity for absorbing large amounts of the refrigerant, often up to 80% of the absorbent dry weight. The complex compounds also exhibit vapor pressure independent of the refrigerant concentration and can be made to absorb and desorb very rapidly. Apparatus using complex compounds to produce cooling are disclosed, for example, in U.S. Pat. Nos. 5,161,389, 5,186,020, and 5,271,239. Improvements in achieving high reaction rates for the complex compounds are achieved by restricting the volumetric expansion of the complex compound formed during the absorption reaction of the gas on the metal salt. The methods and apparatus for achieving such high reaction rates are disclosed in U.S. Pat. Nos. 5,298,231, 5,328,671, 5,384,101 and 5,441,716, the descriptions of which are incorporated herein by reference.\nWhile increased reaction rates have resulted from the aforesaid methods, it has been determined that repeated and relatively long-term absorption and desorption cycling of the complex compounds, particularly those using ammonia as a refrigerant, leads to sorbent migration even in the confined reaction chamber. It has also been found that the sorbent migration increases as higher sorption rates are used. Such migration may lead to uneven sorbent densities which in turn cause force imbalances in the heat exchanger structure, often resulting in deformation of the heat transfer surfaces and/or internal structures. As the heat exchanger structure becomes modified or compromised, heat and mass transfer reductions occur as does the sorption rate of the process. As sorbent migration continues, significant losses in performance efficiency are realized as is the possibility of failure of the reactor especially where it is exposed to high reaction rate sorptions.\nAlthough improvements in attempts to overcome sorbent migration have been made for metal hydrides, such procedures and techniques have not been found to be suitable for ammoniated complex compounds. In U.S. Pat. No. 4,507,263, there is described micro-immobilization for metal hydride using a sintering process in which a metal hydride powder is embedded in a finely divided metal and the mixture sintered in a furnace at 100-200.degree. C. using hydrogen pressure of 250-300 atmospheres. Although such a process reportedly results in mechanical stability for metal hydrides even after 6,000 cycles, the process is not effective for ammoniated complex compounds which exhibit much larger forces as compared to those experienced with metal hydrides. For example, where ammoniated complex compounds are absorbed and/or desorbed above about 3 moles NH.sub.3 /mole sorbent-hr, the forces exercised on a sintered metal structure are so large as to result in deformation of the structure. Moreover, for most practical applications using complex compound technology, practical life expectancy of the reactors will exceed 6,000 cycles by an order of magnitude."}
-{"text": "Work machine operators can experience significant levels of vibration. Many regulatory bodies have imposed restrictions on the vibration levels that an operator may be exposed to over time. To comply with these restrictions, an operator's time on a particular machine can be limited. Specifically, the operator may be required to cease operation of the machine once he has experienced a certain vibration level for a predetermined period of time. Alternatively, an active vibration management system may be employed in an attempt to reduce the average vibration level experienced by the operator and, therefore, prolong his allowed time on the machine.\nVarious systems have been proposed for actively reducing vibrations in a machine. Many of these systems involve sensing of vibrations produced in the machine and reducing the vibrations transferred from a vibration source to the frame of the machine. For example, U.S. Pat. No. 6,644,590 to Terpay et al. (\u201cthe '590 patent\u201d), which issued on Nov. 11, 2003, describes an active system and method for reducing vibrations generated by a gearbox in a rotary wing aircraft. In this system, an active mount is connected between the gearbox and the airframe using hydraulic actuators to suspend the airframe from the gearbox. Based on output signals from various vibration sensors, hydraulic fluid may be supplied to the actuators to move the gearbox relative to the airframe. This motion may be controlled to minimize the transfer of vibrations from the gearbox to the frame.\nWhile the system of the '590 patent may help reduce the vibrations transferred to certain machine components, the system has several shortcomings. For example, the system of the '590 patent cannot monitor or track average vibration levels experienced by an operator or component. Further, the system includes no predictive capability for determining the vibration response of a system to various operator inputs. In addition, the system does not include the capability of adjusting the response of a machine component to reduce the amount of vibration produced. Therefore, the system of the '590 patent may be unsuitable as a means for ensuring that an operator of a work machine does not experience a certain vibration level for greater than a permissible length of time.\nThe present disclosure is directed to overcoming one or more of the problems associated with the prior art active vibration reduction systems."}
-{"text": "It is well known that ambient illumination, that is light originating from sources external to the display device, is reflected to the observer from various optical interfaces of the device and thus reduces the image contrast by increasing the apparent brightness of the dark image areas. Under conditions of high ambient illumination, the image contrast is severly degraded. In addition, a part of the light emitted by the luminescent material of the device also undergoes undesired reflections, producing a further degradation of contrast and of resolution. When the luminescent material consists of a layer of phosphor material in the form of small powder particles, scattering of the emitted light also occurs, further degrading resolution.\nVarious means for overcoming these problems have been proposed. These include the use of various filters including polarizing, neutral density and restricted angle or multi-apertured opaque filters. Other methods include the incorporation of a dark material into the glass of the tube face, or a black dye in the phosphor dielectric layer of the display device. All of the methods have the common disadvantage that the emitted light as well as the reflected ambient light intensity is reduced, with the result that the improvement is contrast ratio is less than desired because the emitted light intensity is a factor upon which the contrast ratio depends.\nThe remarkable reflection-reducing properties of inhomogeneous films were recognized as early as 1880 by Lord Rayleigh (Proc. Lond. Math. Soc. 11, 51, 1880); the properties of such films have been extensively reviewed in a recent series of articles by Jacobsson (Progr. in Optics 5, 247, 1965; Arkiv Fysik 31, 191, 1966; Physics of Thin Films 8, 51, 1975). According to Jacobsson, experimental studies to date have been mainly devoted to transparent inhomogeneous films composed of graded mixtures of two nonabsorbing materials such as ZnS--Na.sub.3 AlF.sub.6, ZnS--CeF.sub.3, CeO.sub.2 --CeF.sub.3, and CeO.sub.2 --MgF.sub.2. These films were found to be durable and of good optical quality. A high index mixture of Ge--ZnS has been produced for application in the infrared wavelength region but were found to be relatively soft and sensitive to moisture and inferior to Ge--MgF.sub.2 films. KBr--Au films were found to have a very low absorption index, with k = 0.01 even at a concentration of gold of 0.16 parts by volume of gold. By contrast, an absorption index of 1.0 was found for a Ge--Au mixture containing 0.1 parts by volume of gold. Ge--In films were also found to have relatively high absorption. Due to the low solubility of In in Ge, the In was expected to remain a separate phase in the form of more or less spherical inclusions.\nAn inhomogeneous Ge--Si.sub.x O.sub.y film was shown by Jacobsson (1965) and also Olsen and Brown (Res./Develop. 16, 52, 1965) to lower the reflectance of a Ge surface to that of a surface of Si.sub.x O.sub.y (refractive index 1.62). Even lower reflectance was obtained with Ge--MgF.sub.2 films, although the transmittance was higher than expected (Jacobsson and Martensson, App. Optics, 5, 29, 1966). One of the first applications of inhomogeneous films as an antireflection coating was described by Nadeau and Hilburn in Canadian Pat. No. 418,289 (1944), and U.S. Pat. No. 2,331,716 (Oct. 12, 1944), in which a plastic layer of polystyrene or urea-formaldehyde resin having a high refractive index is diffused into the surface of an article and overcoated with a second plastic of low refractive index such as cellulose caproate or ethylcellulose. An important commercial application of inhomogeneous films as a low reflectance, absorbing coating on sunglasses was described by Anders in U.S. Pat. No. 3,042,542 (German Pat. No. 1,075,808; 1960). The inhomogeneous films described by Anders consisted of a mixture of low refractive index material, CeF.sub.4, ThF.sub.4, MgF.sub.2, or SiO.sub.2, and a metal, Ni, Fe, Mn, or Cr, or lower oxide of Nb, Ta, or Ti.\nRecently, Steele has proposed in U.S. Pat. No. 3,560,784 the use of a dark dielectric layer consisting of SiO.sub.2 with a tapered concentration of codeposited aluminum applied to the rear side of a light transmissive phosphor layer to serve as a light absorbing layer. The tapered concentration of aluminum results in a continuous variation of the index of refraction through the layer, and such layer comprises an optically inhomogeneous film. Steele claimed novelty for a high contrast cathode ray tube utilizing this construction in which the refractive index of the silicon oxide was substantially equal to that of the phosphor. Phospors suitable for use with the inhomogeneous film of Steele were not otherwise identified. The same objective was the object of an earlier patent of Coltman (U.S. Pat. No. 2,616,057) in which the light absorbing layer was described as lampblack or the black deposits produced by evaporating metals such as aluminum or antimony under poor vacuum conditions.\nUp to the present, the deposition of tapered inhomogeneous films such as in the Steele patent has required the evaporation of two different materials, with the rate of evaporation of each varied as a function of time. Also, it is usually desired that the initial portion of the deposit consist of one component only with the end portion consisting of the second different material only. Steele shows the initial and end materials to be SiO.sub.2 and aluminum, respectively. These requirements pose severe technical difficulties and to achieve reproducible results, elaborate monitoring and control equipment is required so that despite the superior performance offered by inhomogeneous films as compared to homogeneous films, very limited commerical application has been made of inhomogeneous films.\nOsterberg (J. Opt. Soc. Am. 48, 513, 1958) has shown that transmitted waves cannot suffer loss of energy by reflection as they traverse nonabsorbing, inhomogeneous media in which the optical properites have no discontinuities. This result is strictly true only when the medium is infinite in extent. For practical applications, film thicknesses used are of the order of the wavelength of light so that interference due to reflection at the boundaries occurs. The width of the reflectance minimum has been found, however, to be greater than can be achieved with homogeneous films. It also has been shown by Osterberg that inhomogeneous absorbing media similarly cannot exhibit reflectance when the optical properties are continuous. In this case, the medium need not be infinite in extent. Anders (Dunne Schichten fur die Optik, Wissenschafftliche Verlagsgesellschaft mbH, Stuttgart, 1965, English translation as Thin Films in Optics, The Focal Press, London, 1967) has observed that a film thickness of only one wavelength is sufficient for essentially complete absorption in an absorbing inhomogeneous film. This property is basic to the dark dielectric layer described by Steele in U.S. Pat. No. 3,560,784 (1971) since the tapered concentration of aluminum results in an absorbing inhomogeneous film. The deposition of such film entails, however, the technical difficulties previously described, including the deposition of two different materials from two sources."}
-{"text": "1. Field\nThe invention pertains to enhancing the quality of recorded service data, such as data recorded on service tickets, in a data center or call center.\n2. Description of the Related Art\nService delivery centers are large, complex and dynamic ecosystems, which engage hundreds of thousands of experts globally to manage thousands of processes supporting thousands of IT systems with hundreds of configurations. While operations at service delivery centers are typically associated with back-end processes, its efficiency affects quality at front-end (e.g., client experience and satisfaction).\nMultiple ticketing systems, data stores and warehouses trace the operations in service delivery centers. They capture practices of Subject Matter Experts (SMEs), who are typically System Administrators (SAs), and changes in the IT infrastructure (e.g. server decommissioning). These ticketing systems, and enterprise-level warehouses are only reliable as their sources, whether human-driven (tickets submitted by SAs) or system-driven (automated updates of server registries).\nAll too often, there is poor quality of captured data when managing a data center or call center. Administrators are time pressured to achieve high throughput and problem resolution, and no incentive exists for quality of records and logs when capturing and describing problems and resolutions. Low quality of such data leads to inefficiencies in operations (e.g. incomplete tickets slow down the problem resolution process), or leads business analytics to reach wrong or suboptimal conclusions. Frequently, data records such as tickets are blank with insufficient data, and as such are unusable.\nMoreover, low quality of data affects the business decisions (e.g. leading to poor business insights when identifying opportunities for new service offerings, such as \u201cshow me the low utilization servers across the banking sector\u201d). Business insights and problem resolution processes require careful quality assessment to build credibility with stakeholders and efficiently resolve problem tickets. Moreover in such volatile environments, quality of operations and business insights will vary depending on the corresponding data source.\nPlanning activities also depend on good quality data. Take for example server consolidation, where old servers or underutilized servers are migrated into virtual environments with newer hardware. Being able to understand the configuration information such as number of CPUs, speed, memory, operating system and software configured as well as resource information such as network bandwidth, disk and CPU utilization are all key to be able to prepare a plan that maps to proper sized servers. Bad quality data could easily derail a plan from improper source selection to bad target allocations.\nAccumulated problem resolution records contain tremendous source of information about the managed system, its efficiencies and weaknesses, and in addition to analytics, it is a valuable source for knowledge transfer and learning in attempt to train new administrators. The record data are also used for reporting and report generation in billing and service level agreement (SLA) measurements.\nAccurate records of services provided are valuable for a number of business aspects. These include planning of future system improvements, automating problem resolution, optimization of tasks, and awarding the best administrators and skill development. It would be desirable to have a way to improve capturing of incident and problem description and resolution in a data or call center."}
-{"text": "1. Field of the Invention\nThe invention relates to an arrangement for microwave transmission between wave guide regions having different internal gas pressures and/or different fill-gas compositions, that is to say, for coupling or outcoupling microwaves of such a wave guide region into another region.\n2. Description of the Related Art\nIn German Patent Application No. DE-OS 36 22 614 which corresponds to commonly owned U.S. Pat. No. 4,877,642, is disclosed a method of manufacturing electrically conductive moulded bodies by a plasma-activated chemical deposition from a gaseous phase. With such methods the coupling of high-power microwaves is effected through a hermetically sealed insulating microwave aperture of dielectric material in a microwave resonator used as a reaction chamber, in which a plasma is formed and electrically conductive layers are chemically deposited. During this process the problem arises that an electrically conductive film generally covers the surface of the microwave aperture arranged at the coupling place, that is, its inside surface facing the reaction chamber, as a result of which the coupling is stopped. This problem is solved according to DE-OS 36 22 614 either by having the inside of the microwave aperture rinsed by an inertial gas, or selecting for the microwave aperture a dielectric material which is kept free from growth of electrically conductive film as a result of an etching reaction with one of its reaction partners.\nA cognate problem occurs when high-power microwaves from gyrotrons are outcoupled during transition from high-vacuum to air. With microwave powers of the order of 0.1 to 1 MW the thermal load of the known materials used for microwave apertures becomes too large, as a result of which the output power is restricted. With maximum power levels of 0.3 MW one manages by enlarging the wave guide and additionally cooling the aperture consisting of, for example, Al.sub.2 O.sub.3.\nEvacuation of a wave guide through non-radiating or non-coupling slots is known from British Patent Specification No. GB-PS 644,749."}
-{"text": "The rapid rise of health care costs has become an important issue in modern society. To help reduce the costs, professional care givers have begun to seek alternatives, one of which is home health care services. These services not only tend to reduce costs, but are also preferred by the patient wishing to remain in his familiar environment. Among the many types of services provided are: respiratory care, rehabilitation therapy, cardiac monitoring procedures, and infusion therapy.\nInfusion therapy involves IV administration of drugs. Making this therapy safe and convenient for a home situation allows a great number of patients who would otherwise be hospitalized to remain at home and still receive medication. Currently, over 300,000 patients annually use a home infusion therapy delivery system. Typically, patients include the elderly with chronic diseases like cancer, patients with either Crohns disease, HIV or other immune system disorders, and patients suffering from chronic pain. Many of these patients require infusion treatment over a long duration such as months or even years.\nOne characteristic of home IV drug therapy, in contrast to hospital administered therapy, is that a nurse is not always present or readily available. To provide safe and effective treatment, home infusion therapy usually requires that the patient himself, or other non-professional caregiver, such as a relative, administer IV fluids. Special training is required because many home care patients on IV therapy require multiple drugs or multiple doses of the same drug each day. The average nursing visit to a home infusion therapy patient is typically about 90 minutes including commuting time. The typical patient gets between 1 and 4 nursing visits per week, but has to take IV medications daily. Since the cost of daily care by a nurse is not usually covered by most insurers, the cost of attention by a nurse is most economically applied in training the patent or other amateur caregiver and in monitoring the therapy program.\nIn the home care situation non-compliance, over-medication or under-compliance with the IV therapy protocol is a serious issue and quite prevalent. For instance, non-compliance (not taking a medication) or under compliance (taking fewer or smaller dosages than prescribed) occurs in up to approximately one-third to one-half of elderly home therapy patients. Typical compliance related problems include forgetting to follow the specified procedure for administration of the IV medication, forgetting to turn on the various devices used to administer the IV medication and forgetting to turn off a medical device which then delivers too much medication (over-medication). Reasons for compliance related problems are varied and include poor communication, confusion or forgetfulness regarding the procedures and/or equipment, or even attempts to avoid the adverse side effects of IV medications and fluids. Misapplication of the home IV therapy protocol can have serious ramifications resulting in greatly increased home health care nursing expenses, re-hospitalization, and reduction in health status of the patient. Thus, there is a strong need for improved monitoring of patient compliance with the health care program. Benefits of such improved monitoring and compliance include, but are not limited to, improved health at a lower cost, while still remaining in the preferred home environment.\nTo properly monitor compliance with an IV therapy protocol, a device may be provided for monitoring the flow of IV medications and fluids. The IV fluids for a single patient are likely to come from several different sources or systems including IV pumps, IV fluid controllers, gravity drips, syringes, and other devices.\nA typical gravity powered IV may be as simple as an IV bag hanging on a pole in which a nurse or care giver manually adjusts a valve to limit the flow rate, but not control it accurately, or it may use an electronic controller which optically counts the drops of fluid as they pass an optical sensor and then adjusts the flow rate accordingly. However, optical drop counting sensors only provides an indication that the fluid is flowing past the sensor when in a vertical orientation such as hanging from an IV pole. Thus the patient and IV delivery equipment must remain relatively stationary during the administration of the medication or fluid. Optical drop counters also function poorly at higher flow rates and higher line pressures, such as when a syringe is used, because the fluid moving past the drop counter tends to become a continuous stream rather than remaining discrete drops. Therefore, the optical drop counter technique cannot be adapted for use with all fluid sources.\nAn alternative to an optical drop counting sensor, or as a stand-alone measuring device, is a single point pressure transducer to measure the fluid pressure in the IV tubing at a selected point of measurement. This type of sensor is common in IV pumps and is used to indicate that the pump is generating a static pressure head and, correspondingly, causing fluid flow or backpressure in the event of an occlusion in the IV line. This type of sensor only determines line pressure at the selected point, and is only useful in monitoring the pressure caused by the IV pumping device and the related backpressure caused by moving fluids into the patient's body. However, this type of single-point pressure sensor is useful in many IV delivery systems to determine if fluid pressures are at correct levels, and to detect changes in fluid pressure which are indicative of an occluded or collapsed vein. Often, when a certain threshold pressure is detected in a device using this type of sensor, an alarm is sounded to warn of a flow problem. This type of device measures changes in the static line pressure of a fluid line, but is unable to determine if a patient is following proper IV drug administration procedures and cannot differentiate between changes in pressure due to fluid flow versus some other cause, such as an occlusion in which there is actually no fluid flow.\nIncreased backpressure in an IV fluid line causes problems, and, as described above, many IV fluid delivery systems use a sensor to determine when high backpressure develops, i.e. , greater than a few inches of water, for instance when an infiltration of tissue occurs or the tubing becomes occluded. Upon the detection of a significant backpressure, the device sounds an alarm and may function to automatically discontinue the delivery of the IV medication and fluids. Therefore, it is important that any device used to monitor whether or not fluid is flowing does not cause a substantial increase in backpressure or a false occlusion alarm might be triggered.\nOther alternatives use indirect methods to monitor the flow of IV fluids. For instance, the speed and number of rotations in a pump mechanism may be monitored to indirectly determine when fluid flow is occuring. This is useful for flows caused by an IV pump, but is of no value to patients who also receive gravity drips or fluids via syringe. Since nearly all infusion therapy patients must perform venous access device maintenance procedures, such as a heparin flush via syringe to maintain the patency of their IV lines, this pump rotation technique is not of value for monitoring all infusions.\nThe time usage for an IV delivery system may be recorded to prepare bills to patients. Typically, the information is printed or stored in an electronic memory device such as the electronic controls for the drop counter or IV pump. The information may be used to determine which of several patients are using the IV system being monitored, it may be used to coordinate several IV delivery systems with a centrally managed pump, or it may be used to facilitate billing and reimbursement. Unfortunately, none of these systems accommodate tracking of fluid delivered from a variety of sources such as to a patient who receives syringes, gravity drips, and IV pump infusions. The present invention provides an improved flow indicator switch, which overcomes the above-mentioned limitations of the prior art."}
-{"text": "The present invention is directed to a method of optimizing the steering assistance of a motorized vehicle, using angle sensors instead of a torque detector. The method of the present invention also provides an improved steering power in case of failure. The present invention also encompasses a vehicle comprising two angle sensors used to optimize the steering assistance.\nFor utility vehicles, a steering assistance is necessary. It is usually provided through a torsion bar, which opens a hydraulic valve, according to the torque applied by the driver to the steering wheel. In case of failure of the hydraulic pump, or another part of the steering system, the effort to steer the steered axle considerably increases. In case such a failure occurs on an heavy truck, the driver becomes unable to steer the steerable wheels. It is therefore necessary to provide a backup steering system, which allows at least partial steering power. A back up steering system usually requires a second torsion bar, which is costly, heavy and space consuming. It is sometime not possible to implement such a second torsion bar on the steering column. DE102004049038 describes the use of two angle sensors to record the data resulting from the torsion of the torsion bar. However, DE102004049038 is not directed to backup steering systems.\nIt is therefore desirable to provide a method of optimizing the steering assistance with a costly efficient and space saving solution.\nThe steering system of an aspect of the present invention comprises one torsion bar and two angle sensors. The first angle sensor is positioned upstream the torsion bar and the second angle sensor is positioned downstream the torsion bar, in such a way that the torsion angle of the torsion bar can be monitored by the means of the two angle sensors. The portion of the steering column which is upstream the torsion bar comprises all the mechanical elements between the steering wheel and the part just above the torsion bar. It encompasses for example the upper shaft, the lower shaft, with inner shaft and outer shaft, a steering wheel adjustment device. The portion of the steering column which is downstream the torsion bar encompasses all the elements between the torsion bar and the steered wheels. This part comprises for example the drop arm, ball joints, drag link, the upper steering arm, the track rod. In case of twin steered axles, the portion which is downstream the torsion bar also encompasses the elements involved in the steering of the second steered axle. In particular, the second steering pump, the steering actuator of the second steered axle, and the secondary steering rod are downstream the torsion bar.\nIn a first embodiment, the angle sensors are used to detect an abnormal increase of angle between the first and the second angle sensor.\nThe method of the present invention comprises the steps of\na) Monitoring the steering angle of the steering wheel, by the means of a first angle sensor;\nb) Monitoring the steering angle of the steered wheels, by the means of a second angle sensor;\nc) Comparing the difference between the steering angle of the steering wheel, monitored is step a), and the steering angle of the steered wheels, monitored in step b), with a first reference value and/or comparing the steering angle of the steered wheels, monitored in step b), with a second reference value;\nd) Detecting whether the difference between the steering angle of the steering wheel monitored in step a) and the steering angle of the steered wheels monitored in step b) reaches the first reference value of step c) and/or whether;\ne) If the difference between the steering angle of the steeling wheel monitored in step a) and the steering angle of the steered wheels monitored in step b) reaches the first reference value of step c) and/or the steering angle of the steered wheels, monitored in step b) differs from the second reference value of step c), then activating a failure mode.\nIn step a), the angle to which the driver steers the steering wheel is determined by the means of the first angle sensor, positioned upstream the torsion bar. Each angle of rotation of the steering wheel may be associated or not associated to a theoretical angle of rotation of the steered wheels. The theoretical angle of rotation of the steered wheel is the angle expected for a given steering angle of the steering wheel. It may be for example a linear function of the steering angle of the steering wheel. Alternatively, the theoretical angle of the steered wheels may be a non-linear function of the angle of rotation of the steering wheel. The first angle sensor is preferably an angle sensor already present on the vehicle and involved in other functions. For example, the first angle sensor may be the angle sensor already used for the ESP functions.\nIn step b), the effective steering angle of the steerable wheels is determined by the means of a second angle sensor, positioned downstream the torsion bar. This second angle sensor is preferably positioned close to the torsion bar, on the output shaft of the steering gear, in order to provide a direct measurement. However, the second angle sensor may be positioned anywhere else downstream the torsion bar. In case of twin steered axles, the second angle sensor is preferably positioned on the first steered axle. The second angle sensor is preferably an angle sensor already present in the vehicle and involved in other functions. Indeed, an angle sensor may already be present for the steering management of the second steered axle. In this case, there is no need for additional specific sensors.\nStep a) is concomitant with step b). This means that the steering angle of the steering wheel, is determined in step a) at the same time the steering angle of the steered wheels is determined in step b). Monitoring the steering angles in steps a) and b), or the difference of angles, has to be understood as repeating the operation of determining the steering angles, either permanently or as soon as one of the steering angles is modified. Permanently determining the steering angles means that a regular measurement is performed, for example at a predetermined frequency. Preferably, the steering angle is determined each few milliseconds, most preferably between 1 and 10 milliseconds.\nIn step c), the difference between the steering angle of the steering wheel and the steering angle of the steered wheels is monitored and compared to a predetermined value, which is a first reference value, or a warning threshold value, under which should remain the difference of steering angles. If a theoretical value is associated to the steering angle of the steering wheel in step a), the effective steering angle of the steered wheels, measured in step b), may also be monitored and compared to this theoretical value, which is a second reference value. Under normal conditions, the effective steering angle of the steered wheels should correspond to the second reference value. Also, under normal conditions, the difference of the steering angles determined in steps a and b) should remain under the first reference value. Under these circumstances, it is considered that the suitable steering assistance is delivered, allowing effective steering of the steered wheels. No additional steering power is triggered.\nIn step d), it is identified that the difference of the steering angles, reaches the first reference value or the effective steering angle of the steering wheels departs from the second reference value. Under these conditions, it is considered that the steering system is in fault and step e) is initiated. Alternatively, step e) may be initiated if the two conditions of step b) are reached. In this case, step e) is initiated only when the difference of the steering angles reaches the first reference value and the effective steering angle of the steering wheels departs from the second reference value.\nStep e) triggers a failure mode, wherein additional power steering is delivered to compensate the efforts of the driver. The failure mode may be the activation of an auxiliary steering power. In case of more than one steered axle, the failure mode may be a special mode of the steering system of the second steered axle. For example, under failure mode, the steering system of the second steered axle may be activated in a way to provide an oversteering of the second steered axle. The failure mode may encompass any other action which aims at improving the steering assistance."}
-{"text": "An electronic device includes a connector into which a cable for wired communication is inserted and a slot into which a removable storage medium is inserted in order to exchange information with another device. Therefore, the housing of an electronic device includes a hole for exposing an electronic device terminal such as a connector or slot to the outside of the housing. Hereinafter, in the present specification, a hole for exposing an electronic device terminal to the outside of the housing is referred to as a \u201csocket\u201d. Generally, the socket of an electronic device is covered and protected with a cover.\nRecent electronic devices are required to be waterproof and need to prevent water from entering through the sockets. Therefore, packings are provided on covers that cover the sockets to protect the electronic device terminals. Structures for preventing water from entering through the sockets are classified into a structure called a longitudinal compression type in which a packing is placed on the front side of the opening of the socket and a structure called a transverse compression type in which a packing is put in the socket and brought into contact with the wall. The longitudinal compression type needs to secure a packing margin around the opening of the socket, and thus hinders miniaturization of the electronic device, so the adoption of the transverse compression type is progressing.\nPatent Literature 1 discloses a transverse compression type waterproof structure."}
-{"text": "For the ignition of briquets and charcoal for grills up till now ignescent fluid has been the dominating and sole accepted ignition aid for producing in an acceptably short time embers for broiling. Among the drawbacks of ignescent fluid are the hazards of the ignition procedure. At times, ignescent fluid has been confused with other fluids and caused severe burns in children and in some known instances children have been poisoned by drinking the fluid. In addition, the ignescent fluid is bulky and generally difficult to bring along. It sometimes also imparts obtrusive flavours to the food being broiled. The use of ignescent fluid is also expensive.\nTo light a fire in fireplaces, furnaces, and suchlike, one normally uses newspaper leaves and the like, in conjunction with wood chips. This is a time-consuming method. Ignition aids known as `fire lighters` may also be used. A method for producing fire lighting aids was described in SE-A-No. 41 897, in 1914. According to this method, paper, sulphite or sulphate pulp, is impregnated with a combustible substance which is either liquid or solid, such as resin, resin dissolved in some combustible substance such as spirits, turpentine, raw or refined petroleum, tar, or some other suitable substance. After being impregnated, the paper or the pulp is rolled onto spindles, and fire lighting aids then prepared from the strips, whether wet or dry, the final product being in the form of small reels. According to SE-A-No. 96 174 fire lighters are produced from lumbering or wood mill debris, which is cut into chips, defibrated, mixed with water to achieve a suitable consistency and lastly formed into a plate, which is dewatered by pressing and then dried. This plate is dipped in molten paraffin, stearin, or tallow or a mixture of these at a temperature of 80.degree.-100.degree. C. After drying, the plate is cut into pieces of a certain width and length. Before being impregnated, the plate is provided with grooves, to facilitate the cutting of the plate into small square blocks.\nA drawback which is common to these and other known fire lighting aids is that the area of combustion is small, the product thus having to be ignited at a very small area. Therefore, it is not at all uncommon to fail at the ignition of these products, even if the burning time may be long. In addition, the positioning of the lighter is critical, for instance when lighting a fire on a grill, since the lighter, being very small, may easily fall down through the grid.\nAnother known lighting aid consists of cubes of a brittle material which easily crumbles and has a strong odour, so that the product must be carefully packed and gently handled.\nParaffin impregnated cellulose pulp is a better lighting aid. The area of combustion of this product in relation to its volume is greater, and hence the product burns more intensely and over a larger area. Even though its burning time is shorter than that of a more compact product of the same volume, the fire or the bed of briquettes or coal to be ignited is lit more effectively and more safely. Another desirable property of the lighter is that it is free of tackiness. Nor should it crumble when broken, as is the case if not all paraffin has become absorbed into the pulp. At the same time it must be water-repellent and inflammable. These demands have caused considerable manufacturing problems."}
-{"text": "This application claims the priority of German Patent Application No. 101 54 669.6, filed Nov. 7, 2001, the disclosure of which is expressly incorporated by reference herein.\nThe present invention relates to an internal combustion engine having at least two cylinder banks and more particularly, to an internal combustion engine whose cylinder heads are sealed by cylinder head covers, wherein to ventilate the crankcase from the so-called blow-by gases, ventilation lines are connected to the cylinder head covers and communicate with a negative pressure source, e.g., an intake pipe, and on the inside of the cylinder head cover means are provided for pre-separating the oil from the blow-by gases.\nU.S. Pat. No. 3,908,617 discloses a device for crankcase ventilation of an internal combustion engine with two cylinder banks in which ventilation lines mounted above the cylinder head housing or the cylinder head cover remove the blow-by gases located in the crankcase volume and return them to the intake system of the internal combustion engine in a closed circuit. In addition, sheet metal guide elements are mounted on the inside of the cylinder head cover. The blow-by gases flow past these guide elements and a portion of the oil carried along by the blow-by gases is deposited thereon."}
-{"text": "The present invention relates to temperature-producing conductive-resistive medium and to a method of producing a variety of articles therefrom.\nThere have been many attempts to produce electrically-conductive coatings such as paints. Generally, there are two types of electrically-conductive coatings. The first is a low resistivity, high conductivity paint that contains a pigmentation of metal particles while the second is a high resistivity, low conductivity paint that is formed from compositions containing carbon or graphite.\nLow resistivity paints have traditionally been used to provide coatings having a high conductivity for connecting conductors that require a superior electrical bond with a minimum resistance. Generally, low resistivity paints cannot be applied to materials in order to produce temperature adjustable heating elements because the low resistivity paint requires a high volume of current to generate a reasonable output of heat. In contrast, the resistivity of traditional highly resistive paints is often so high that a relatively high voltage drop is required in order to generate sufficient heat. As a result, the use of high resistivity paints usually sacrifices safety. Furthermore, when either of the above-identified traditional conductive paints are applied to various substrates, cracks and flaking of the paint often develop over a period of time. This causes a breakdown in the temperature adjustable property of the article.\nIt is therefore an object of the present invention to provide a method and apparatus for generating an electrical resistance temperature adjustable substance for application to a variety of substrates in order to provide temperature controllable properties.\nIt is another object of the present invention to provide a method and apparatus for generating an electrical resistance temperature adjustable substance for application to a variety of materials wherein the electrical resistance temperature adjustable substance does not inhibit the inherent flexibility of the substrate to which it is applied.\nOther and further objects will be made known to the artisan as a result of the present disclosure and it is intended to include all such objects which are realized as a result of the disclosed invention."}
-{"text": "From DE 1 103 216 a device for distributing cut tobacco to cigarette-making machines is known, wherein the cut tobacco is fed from a conveyor onto a rotary table from which the tobacco is drawn by stationary sucking pipes spaced at the periphery of a table constituting a distributing element, the cut tobacco fed from the conveyor falling onto a cone located centrally relative to the rotary table. The cut tobacco slides down along the cone onto the rotary table gravitationally and then it is transported due to the centrifugal force as a layer towards the periphery of the table, from where it is sucked by vertical pipes to deliver the cut tobacco to the cigarette-making machines.\nDE 198 23 873 presents a similarly operating device for feeding cut tobacco to many machines. The cut tobacco is fed via a vertical channel onto a bowl performing a composed, rotary and circulating, motion. The sucking channels, picking up the cut tobacco from the uniformly formed layer, are arranged vertically within the bowl cover at the bowl periphery.\nIn GB 959 343 a device is described in which the cut tobacco is fed, as previously, from above onto a rotary distribution disk and is directed by the centrifugal force towards receiving channels arranged radially in the side wall of the distribution chamber.\nIn a slightly different arrangement, known from DE 300 90 000, cut tobacco is fed through a charging hopper onto a linear vibrational conveyor. The vibrational conveyor transfers the fed cut tobacco to a place above which sucking pipes are situated. The cut tobacco is transported in the form of a layer and the sucking pipes are arranged vertically just above the surface of this layer.\nUsually the bottom of the distribution chamber is flat or has the shape of a bowl and it is a surface of revolution and posses a centrally located rotational cone.\nThe process of feeding the cut tobacco to the cigarette-making machines is discontinuous, the result of which is that the more receiving channels are connected, the more frequent changes of the flow rate of the tobacco through the distributing device will occur. The discontinuity of the feeding process results from the fact that after filling the cut tobacco container located within the machine, the feeding is stopped until the amount of the cut tobacco in the container drops below a certain predefined level, afterwards the feeding is started again. Devices for distributing cut tobacco, employed in the tobacco industry, usually feed a lot of cigarette-making machines. Every change in a total throughput of the receiving channels will result, as a consequence, in a change of the efficiency of the conveyor feeding the distributing device.\nAll the solutions presented above relate to devices for distributing cut tobacco to cigarette-making machines using gravitational feeding, usually in the form of a feeding channel and a couple of pneumatic receiving channels transferring the cut tobacco to the cigarette-making machines, the receiving channels being connected to the distributing chamber or being located at the periphery of the distributing element for uniform distributing the cut tobacco into the inlets of the receiving channels. For proper operation of all the above devices it is necessary to collect some amount of the cut tobacco in the distribution chamber, which is transferred to the space from which it is received by the receiving channels. During transferring the layer of the cut tobacco gains its optimal thickness in order to ensure repeatable conditions of receiving the cut tobacco by the receiving channels. Therefore the receiving channels are distant from the feeding channel. In each of the devices in the case of temporary stopping the process of feeding the cigarette-making machines, the amount of the cut tobacco, which has been already delivered to the distributing device but has not been yet received, is an excess of the cut tobacco present in the device relative to the amount necessary for its operation. The cut tobacco tends to agglomerate, i.e., to create bundles, the effect of the agglomeration being particularly strong if the cut tobacco is stored in a high layer, as in the vertical channel feeding the distributing device.\nIf the process of receiving the cut tobacco by the cigarette-making machines, connected to a single distributing device, is stopped, one must stop the conveyor feeding the device, which was operating with a rate adjusted for feeding all the cigarette-making machines. However, due to inertia of the system, the distribution chamber will be filled anyway as well as, partially or fully, then vertical feeding channel. Restarting the device after a longer downtime may occur difficult, since the bulk density of the cut tobacco collected and stored under a pressure within the feeding channel increases and it is significantly more difficult to form a uniform layer of the cut tobacco and to suck the agglomerated tobacco through the receiving channels. Sometimes, in order to restart the feeding system the agglomerated tobacco must be removed from the lower portion of the feeding channel and partially from the distribution chamber.\nIf a couple of receiving channels will be shut off simultaneously, i.e., in the case of a rapid drop of the received amount of the cut tobacco, an excess of the cut tobacco will arise within the distribution chamber. The efficiency of the conveyor feeding the distributing device will be adjusted to the throughput of the cigarette-making machines that are still working, and the excess of the collected cut tobacco will be used by those machines, however if the excess is relatively large, disturbances in the receiving process may arise.\nFrequently, cigarette manufacturers must face the task of producing short series of new cigarette brands. Large distributing devices with rotary tables or vibrational conveyors are expensive and there is no economical justification for using them in the case of frequent changes of the brand of tobacco fed to one or two cigarette-making machines."}
-{"text": "1. Field of the Invention\nThe present invention generally relates to methods and systems for inspection of an entire wafer surface using multiple channels. Certain embodiments relate to detecting light scattered from different portions of the entire wafer surface using different detection channels.\n2. Description of the Related Art\nFabricating semiconductor devices such as logic and memory devices typically includes processing a specimen such as a semiconductor wafer using a number of semiconductor fabrication processes to form various features and multiple levels of the semiconductor devices. For example, lithography is a semiconductor fabrication process that typically involves transferring a pattern to a resist arranged on a semiconductor wafer. Additional examples of semiconductor fabrication processes include, but are not limited to, chemical-mechanical polishing, etch, deposition, and ion implantation. Multiple semiconductor devices may be fabricated in an arrangement on a semiconductor wafer and then separated into individual semiconductor devices.\nWafers may contain defects both in central portions of the wafers as well as in edge portions of the wafers, which includes a relatively narrow region around the periphery of the wafers, and on the outer edge of the wafers. Examples of defects that may be found in the edge portion and on the outer edge of wafers include, but are not limited to, chips, cracks, scratches, marks, particles, and residual chemicals (e.g., resist and slurry). As wafer sizes continue to increase, both wafer and integrated circuit (IC) manufacturers are becoming more concerned about defectivity at or near the wafer edge. The main concerns are that edge defects could fall onto the central part of the wafer thereby causing untraceable yield loss, cross contamination during processing, and/or catastrophic wafer breakage. These yield loss mechanisms are experienced by most wafer and IC manufacturers at one time or another.\nTraditionally, wafer inspection tools are designed to inspect a central portion of the wafers (i.e., a surface area of the wafer on which electrical elements will be formed or a surface area of the wafer opposite that on which electrical elements will be formed). Since these areas of the wafer reflect or scatter relatively small amounts of light, such wafer inspection tools are designed to detect relatively small amounts of light. However, near the outer edge of the wafer, relatively large amounts of light may be reflected or scattered from the wafer due to edge features such as a bevel formed at or near the outer edge. As a result, these large amounts of light will saturate the detectors of traditional wafer inspection systems. Consequently, any output signals generated near or at the edge of wafers by such wafer inspection tools are generally unusable. In some instances, the wafer inspection systems may be designed to block the light from reaching the detectors when inspecting near the edge of the wafer to protect the detectors from damage that may be caused by the relatively high intensity light.\nSome edge inspection systems are being developed to detect defects at or near the outer edge of wafers. Examples of apparatuses for detecting defects along the edge of electronic media such as semiconductor wafers are illustrated in U.S. Patent Application Publication Nos. 2003/0030050 by Choi and 2003/0030795 by Swan et al., which are incorporated by reference as if fully set forth herein. Due to the substantially different reflecting and scattering characteristics of the outer edge of wafers in comparison to the inner portion of the wafer, such edge inspection systems have substantially different configurations than the traditional wafer inspection tools. Therefore, the edge inspection systems are not optimized to, or even able to, detect defects in the central portion of the wafers. Consequently, if wafer or IC manufacturers want to detect defects in both the central and outer portions of wafer (as is usually the case since defects in either portion may result in expensive yield losses and other problems), they will need to purchase two separate tools. Using two different wafer inspection tools instead of just one inspection tool will obviously increase costs in many ways such as increases in clean room real estate and operating costs, increases in tool maintenance costs, and increases due to reduced throughput. However, since a tool that is capable of inspecting both the inner and outer portions of wafers is not currently available, and due to the increasing costs associated with defect-based yield losses, wafer and IC manufacturers may not be able to avoid the costs associated with multiple, different inspection tools.\nAccordingly, it may be advantageous to develop a wafer inspection system that is capable of inspecting substantially an entire surface of wafers including both center and edge portions of the wafers."}
-{"text": "1. Field of the Invention\nThe present invention relates to a channel allocation method of a wireless network and a system thereof. More particularly, the present invention relates to a distributed channel allocation method of a wireless mesh network and a system thereof.\n2. Description of Related Art\nIn recent years, there is a rapid development in the field of wireless broadband access techniques including Wi-Fi (IEEE 802.11 series), WiMAX (IEEE 802.16 series) and 3G, etc. The wireless mesh network (referred to hereinafter as WMN, IEEE 802.11s) is one of the key techniques integrated with the wireless broadband network. The structure of the WMN illustrated in FIG. 1 is a mesh network based on a wireless transmission interface, and the WMN has a similar operation mode to that of an Ad-hoc network. Since the operation of the WMN is based on the wireless transmission interface, it has the advantage of rapid deployment without restriction of the geographical landforms. The WMN is generally applied to a community network, a temporary network of exhibition halls or shopping stalls, networks established in disaster areas or areas having special geographical environments, and so on.\nThe operation of the WMN is based on the wireless transmission interface. Taking the IEEE 802.11a/g for an example, its transmission bandwidth of data is 54 Mbps (mega bytes per second), which is the maximum possible transmission bandwidth. However, influenced by a MAC (media access control) contention, 802.11 headers, 802.11 ACK signals and packet errors, an average applicable bandwidth is usually less than half of the maximum bandwidth.\nFurthermore, the most serious issue lies in that a data transmission rate of a network link layer may be decreased greatly due to signal interference. Two possible interference problems are shown in FIG. 2: (1) interference in the same transmission path, (2) interference in the adjacent transmission paths. Referring to FIG. 2, the signal coverage of a node 3 includes nodes 2, 4 and 9. Similarly, the node 3 is simultaneously in the signal coverage of the nodes 2, 4 and 9. A first transmission path and a second transmission path are paths for data transmission. The first transmission path is taken for an example. When the node 2 and the node 3 are transmitting data, the node 4 may receive signals from the node 3, resulting in the fact that node 4 cannot transmit data to a node 5 provisionally. Therefore, the bandwidth of the first transmission path is reduced, which refers to the so-called interference in the same transmission path.\nOn the other hand, referring to the node 9 on the second transmission path, since the node 9 is in the signal coverage of the node 3, the node 9 may receive signals from the node 3 when the node 2 and the node 3 are transmitting data, resulting in the fact that the node 9 cannot transmit data to a node 8 or a node 10 provisionally. The phenomenon indicating an interference of data transmission through the first transmission path with that through another transmission path (a second transmission path) represents the so-called interference in the adjacent transmission paths. Therefore, many studies are performed on the WMN to learn how to improve an applicable bandwidth of the WMN by advancing a structural design thereof.\nAccording to the IEEE 802.11s WiFi Mesh standard, a plurality of wireless transmission interfaces is allowed to use different non-overlapping channels for transmission, so as to increase the transmission bandwidth. Therefore, some studies have been developed to increase a network flow by applying multi-network interface cards (referred to hereinafter as Multi-NIC). A method of increasing the network flow includes allocating a plurality of NICs on each node, and each of the NICs may employ a different non-overlapping channel to communicate with other nodes. The advantage of this method lies in that it is unnecessary to modify any existing hardware structures. Only is an integral channel allocation method required for assisting the existing hardware structure, and the network flow can be substantially improved.\nA method and a system for assigning channels in a wireless local area network (WLAN) is disclosed in U.S. Publication No. 2006/0072502 A1, in which the WLAN infrastructure mode (i.e. a client to hub communication mode) is provided. A mobile node (referred to hereinafter as MN) in the network is connected to an access point (referred to hereinafter as AP) by means of one hop, and the other end of the AP is connected to a wired network, wherein each AP has at least two applicable channels, and each AP is at least adjacent to another AP.\nEach AP constantly collects the traffic load information and forecasts a possible throughput on each channel. Thereafter, the AP determines an optimal channel for connecting with the MN within the signal coverage of the AP. However, this channel allocation method only takes the optimal channel within one hop between the AP and the MN into account. Therefore, the application of the method is limited.\nMost of the early studies focus on modifying an MAC layer protocol of the network to support a multiple channel transmission. The studies aim to find the optimal channel for transmitting every single packet, so as to avoid the interference. On the other hand, a concept of a Multi-NIC disclosed by V. Bahl et al. and P. H. Hsiao et al. in two articles has drawn attention and discussions recently. One of the articles was authored by V. Bahl, A. Adya, J. Padhye, A. Wolman, entitled \u201cReconsidering the Wireless LAN Platform with Multiple Radios\u201d Workshop on Future Directions in Network Architecture (FDNA-03), while another one was authored by P. H. Hsiao, A. Hwang, H. T. Kung, and D. Vlah, entitled \u201cLoad-Balancing Routing for Wireless Access Networks\u201d Proc. of IEEE Infocom 2001. The method disclosed therein is to install a plurality of the NICs on each node of the Ad-hoc network, and each NIC may dynamically determine a channel for communicating with other nodes. The advantage of this method lies in that it is unnecessary to modify any existing hardware structures. Only is the integral channel allocation method required for assisting the existing hardware structure, and the network flow can be substantially improved. Sequentially, a channel allocation method based on a centralize structure was disclosed by A. Raniwala, K. Gopalan, T. Chiueh, entitled \u201cCentralized channel assignment and routing algorithms for multi-channel wireless mesh networks,\u201d ACM Mobile Computing and Communications Review 8 (2) (2003), which is one of the earliest articles having a formal definition of the channel allocation. In the method, a load-aware channel assignment is performed by an evaluation matrix defined by the authors themselves, the entire network is calculated in overall, and a preferable channel allocation is obtained. Thus, a maximum network flow is then achieved.\nIn recent studies, a channel allocation technique based on a dynamic & distributed structure has been disclosed, wherein channel allocation information is exchanged by using a common channel framework according to the IEEE 802.11s standard. This technique is based on IEEE 802.11 WLAN standard, wherein a plurality of wireless NICs is installed to support a multi-channel transmission. However, the interference still cannot be avoided in the aforementioned techniques."}
-{"text": "Power amplifiers for cellular handsets are optimized for efficiency at, or close to, maximum output power. However, in the field, they may only be called upon to operate near maximum output power for a very small percentage of the time. The rest of the time, they may be operating at back-off output power levels, where their direct current (DC) to radio-frequency (RF) conversion efficiency is very much reduced. This reduced efficiency under practical conditions results in wasted battery power in the handset and, therefore, reduced talk time."}
-{"text": "The present invention relates to ultrahigh molecular weight polyethylene and polypropylene fibers having high tenacity, modulus and toughness values and a process for their production which includes a gel intermediate.\nThe preparation of high strength, high modulus polyethylene fibers by growth from dilute solution has been described by U.S. Pat. No. 4,137,394 to Meihuizen et al. (1979) and pending application Ser. No. 225,288 filed Jan. 15, 1981.\nAlternative methods to the preparation of high strength fibers have been described in various recent publications of P. Smith, A. J. Pennings and their coworkers. German Off. No. 3004699 to Smith et al. (Aug. 21, 1980) describes a process in which polyethylene is first dissolved in a volatile solvent, the solution is spun and cooled to form a gel filament, and finally the gel filament is simultaneously stretched and dried to form the desired fiber.\nUK patent application GB No. 2,051,667 to P. Smith and P. J. Lemstra (Jan. 21, 1981) discloses a process in which a solution of the polymer is spun and the filaments are drawn at a stretch ratio which is related to the polymer molecular weight, at a drawing temperature such that at the draw ratio used the modulus of the filaments is at least 20 GPa. The application notes that to obtain the high modulus values required, drawing must be performed below the melting point of the polyethylene. The drawing temperature is in general at most 135.degree. C.\nKalb and Pennings in Polymer Bulletin, vol. 1, pp. 879-80 (1979) J. Mat. Sci., vol. 15, 2584-90 (1980) and Smook et al. in Polymer Bull., vol. 2, pp. 775-83 (1980) describe a process in which the polyethylene is dissolved in a nonvolatile solvent (paraffin oil) and the solution is cooled to room temperature to form a gel. The gel is cut into pieces, fed to an extruder and spun into a gel filament. The gel filament is extracted with hexane to remove the paraffin oil, vacuum dried and then stretched to form the desired fiber.\nIn the process described by Smook et al. and Kalb and Pennings, the filaments were non-uniform, were of high porosity and could not be stretched continuously to prepare fibers of indefinite length."}
-{"text": "1. Field of the Invention\nThis invention generally relates to an apparatus for recording digital information on a recording medium such as a magnetic disk or tape and reproducing the recorded information and more particularly to a code-error correcting device for correcting code-errors in digital signals used in the apparatus for recording digital information on the recording medium and reproducing the recorded information.\n2. Description of the Related Art\nReferring first to FIG. 8, there is shown the structure of a code of digital signals used in a typical apparatus for recording audio signals on a magnetic tape by means of a rotary head and reproducing the recorded information (that is, what is called an R-DAT (Digital Audio Tape recorder)) therefrom. As shown in this figure, the code includes data (DATA) composed of 28.times.26 symbols, a transverse or horizontal parity code (C.sub.2 PARITY) composed of 28.times.6 symbols and a longitudinal or vertical parity code (C.sub.1 PARITY) composed of 4.times.32 symbols. In the case of Reed Solomon Code (R.S.C), sets of data concerning the parity codes C.sub.1 and C.sub.2 are (32, 28, 5) and (32, 26, 7), respectively. In each of the parentheses, a first, second and third numeral indicates values of the total length of a code, the length of data and a minimum distance between code words, respectively.\nFurther, referring now to FIG. 9, there is shown the format of signals employed when recording the signals having such structure of codes. In this figure, reference characters SYNC indicates a synchronizing signal; ID an identification signal; ADR an address signal; P a block parity signal; DATA data of 28 symbols; and C.sub.1 a C.sub.1 parity code of 4 symbols. That is, signals SYNC, ID, ADR and P are added to data signals. Incidentally, in this case, the block parity signal is given by EQU P=ID.sym.ADR.\nNamely, the signals, of which the format is as shown in FIG. 9, are recorded on the magnetic tape and reproduced therefrom.\nThe above described conventional apparatus can detect address errors to some extent by transmitting the block parity signal indicating ID.sym.ADR together with the signal indicating data DATA. However, the conventional apparatus has a drawback that the capability of detecting the address errors is not sufficient to precisely detect the address error and as a consequence the address errors increase. In this case, data are stored in an erroneous area within a memory in accordance with the erroneous address information because the address information generally determines an area in the memory in which data are to be stored. Conventionally, even when the error cannot be detected by using the longitudinal parity code (C.sub.1), the error can still be corrected if the error is present within the range which can be corrected by using the transverse parity code (C.sub.2). Further, if the error exceeds the capability of detecting the error by using the transverse parity code (C.sub.2), it is necessary to locate the error on the basis of the error information which is generated after the check by using the parity codes C.sub.1 and C.sub.2.\nIn such case, if only the area in the memory is erroneous and the parity code C.sub.1 is correct, there is the inconvenience that in spite of the fact that a sequence of data is erroneous, the error cannot be detected. Thus, to eliminate such inconvenience, it has been proposed that when the parity code C.sub.1 is generated, the address is included as a generating element for the parity code C.sub.1. Such approach has a defect that the capability of correcting error is degraded because the code is not a product code.\nTherefore, it is an object of the present invention to provide a code correcting device of which the capability of correcting error is significantly improved, thereby decreasing the possibility of passing over the error."}
-{"text": "Wireless communication can be used as a means of accessing a communication network. Wireless communication has certain advantages over wired communications for accessing a network. For example, implementing a wireless interface can eliminate a need for a wired infrastructure thereby reducing the cost of building and maintaining network infrastructure. In addition, a wireless network can support added mobility by allowing a wireless device to access the network from various locations or addresses. A wireless interface can comprise at least one transceiver in active communication with another transceiver that is connected to the network.\nVarious types of network configurations can be used to communicate data over the wireless network. For example, a heterogeneous network can be configured to include various types of access nodes such as a macro access node, a micro access node, a pico access node, a femto access node, etc. In a heterogeneous network, a wireless device can be served by an access node having the lowest signal path loss rather than by an access node having the strongest signal strength as in traditional network configurations.\nIn a heterogeneous network, interference can occur at the cell edge of the short range, low power access nodes due to the macro access node. This interference can result in undesirable reduction in coverage and throughput to the wireless devices in communication with the short range access node. A scheduling scheme comprising almost blank subframes (ABS) can be used to create an opportunity for the wireless devices within the cell edge region of a short range access node to receive downlink information without interference from the macro access node. However, ABS subframes can undesirably limit an amount of resources allocated to wireless devices during each frame."}
-{"text": "Bed-type massage devices are generally provided with massaging rollers which can be displaced by a drive device in a bed base and are so constructed that the massaging rollers travel on guide rails provided on both sides inside the bed base, the arrangement being such that when a user lies face upwards on the bed-type massage unit the massaging rollers roll along his or her dorsal region and effect finger-pressure type massage thereof.\nIn recent years there has been research on and development of units in which such bed-type massage devices are made contractible to as compact a size as possible in order to reduce the space needed at times of transport and delivery or during storage.\nFor example, as disclosed in Japanese Laid-open Patent Application No. 59-189848 and Japanese Laid-open Utility Model Application No. 61-54834, there are known bed-type massage devices which are designed to reduce space requirements at times of transport and delivery or during storage by being constructed in such a way that foldable guide rails are provided on both sides inside a bed base and the bed frame itself is formed as an elastic body which is foldable and the massaging rollers can run along the guide rail, whereby the base unit can be folded into halves or thirds when it is not in use.\nHowever, the structure in these conventional bed-type massage devices is such that extension and contraction of the bed base is employed solely for the purposes of transport and delivery or storage and, although the size of the bed base can be reduced in the longitudinal direction, it is impossible to reduce the volume and so the devices still fail to resolve problems in packing, etc.\nFurther, since conventional bed-type massage devices are units designed for the purpose of reduction of size at times of transport and delivery or during storage and are not constructed in a manner permitting adjustment of the extension or contraction in accordance with the height of the user, if a user who is shorter than the length of the bed base uses the base to effect massaging, the massaging rollers move to portions beyond the body of the user, i.e., they move over an unnecessarily large range, which is wasteful both in terms of electric power and of time.\nBy way of a means for resolving this problem, the present Applicant discovered a means whereby a bed base of a bed-type massage device can be extended or contracted in opposed lengthwise directions without being folded and which permits fine adjustment of the extension in accordance with the height of the user. In this case, however, although extension and contraction in opposed lengthwise directions and fine adjustment of the extension are possible, there are problems associated with aspects such as the means for disposing the massaging rollers inside the bed base to match the bed base in different states and the drive means for causing the massaging rollers inside the bed base to move forward and back in a manner such as to match the bed base in different states.\nIt is the object of the present invention to resolve the various problems noted above and provide a bed-type massage device in which a massage unit is not driven when the bed base is contracted but the massage unit can longitudinally travel repeatedly and smoothly over the whole area of the bed base when the bed base has been extended or contracted and fine adjustment of the extension has been made.\nIt is a characteristic of the bed-type massage device of the invention that it comprises a variable bed base which is made freely extendible and contractible in opposed lengthwise directions and in which a first side of a flat second base fits into an opening on one side of a flat first base, at least one retention hole is formed on the left and on the right in the lower surface of the second base, and lock pins which are engageable in these retention holes are so provided that lock mechanisms provided at the tips thereof face the second base from the left and right of the other side of the first base; a pair of guide rails which are laid lying along the longitudinal direction on the left and right of the upper surface of at least the first base of the said bed base; a drive mechanism in which a drive motor is provided at and orthogonal to one end of the second base in the longitudinal direction, worms are respectively coupled with the ends of two drive shafts of this drive motor and the respective ends of a pair of rod-like screwshafts which are respectively provided along both lengthwise sides of the second base are in screw engagement which the respective worms; a massaging unit which is constituted by effecting screw engagement of both sides of the base end of a frame unit with the pair of rod-like screwshafts of the said drive mechanism providing facing running rollers which are capable of moving and running over the said guide rail on both sides of the lower part of the far end of the said frame unit and providing rolling members at a set interval above the base end and the far end of the said frame unit; and a covering cloth that is wound around to cover the entire upper surface of the bed base, and the device is so constructed that the massage unit can repeatedly travel over the bed base and the variable bed base is freely extendible and contractible to permit adjustment to any required length.\nHaving the above construction, the bed-type massage device of the invention brings about the following effects.\nAll that is needed if it is wished to adjust the length of the bed-type massage device of the invention to match the height of a user is to push the first base of the variable bed base in the direction of the second base and lock it in the said position by means of the lock mechanisms, and if it is required to shorten the bed at times such as when the bed is packed for transport or is delivered it can easily be reduced to about half its length by releasing the lock mechanisms and pushing in the first base into which the second base is inserted.\nFurther, when it is wished to extend the bed base after it has been shortened, this can easily be done by simply releasing the lock mechanisms, pulling the first base so that the second base is moved out from it and then, after the bed has been set to the required length, locking it with the lock mechanisms.\nWhen the adjustment of the length of the bed-type massage device has been completed, the user lies face upwards on its variable bed base and then simply actuating the drive motor brings about an agreeable massaging action in which, through the action of the drive mechanism, the massage unit is caused to move and travel on the guide rails, limit switches at the variable bed base's terminal and start ends in the longitudinal direction are actuated and cause repeated displacement and travel of the said massage unit, so causing rotatable rollers of the rolling members to slide while coming into uniform contact with the whole surface of the user's back.\nThe bed-type massage device of the invention will now be described in detail with reference to one embodiment thereof which is shown in the drawings."}
-{"text": "Various applications of fluorescence techniques to analyze biological samples are known to people skilled in the art. In case of electrophoretic techniques proteins or DNA are labeled with a fluorescence probe to visualize their electrophoretic bands in gels or columns. In addition, most biochip applications so far are based on a fluorescence read-out, whereas the specific binding of a fluorescence-labeled target molecule to a probe molecule immobilized on a solid support is monitored. Applications for DNA analysis in the liquid phase include fluorescence hybridization probes like the double-stranded DNA binding dye SybrGreenI or FRET (Fluorescent Resonance Energy Transfer) probes utilizing two fluorescence probes and energy transfer. A very important application for fluorescence techniques in the liquid phase is the quantification of PCR products in real time, the so-called real-time PCR.\nIn all these cases, a fluorescence reading device is needed that provides light of a certain wave length to excite the fluorescence label of the assay and that is able to detect the fluorescence light form said label emitted at a somewhat different wavelength. One major problem of all fluorescence reading devices is the enormous intensity of the excitation light in comparison with the fluorescence light emitted by the dye and therefore, one has to assure that the excitation beam does not hit the detector in order to monitor the fluorescence signals accurately. In other words, the optical path of the excitation light has to be different from the optical path of the fluorescence light, at least partially.\nThe realization of the fluorescence principle is quiet easy, when only one fluorescence probe has to be monitored in the liquid phase of e.g. a capillary. Here, e.g. a white light source together with a set of dichroic mirrors and filters is sufficient to meet the requirements. However, if more than one fluorescence label is present in the sample, a lateral distribution of spots on a solid support or the fluorescence of a microtiter plate has to be monitored, the requirements for the fluorescence reading device are more difficult to fulfill.\nIn principle, there are two different strategies to excite and monitor the fluorescence of a lateral distribution of sites. The first strategy is to scan the lateral distribution of sites, whereby the individual sites are successively analyzed one at a time. The second strategy is to illuminate the whole distribution of sites simultaneously and to image the corresponding fluorescence e.g. on a CCD chip. The scanning strategy has the obvious drawback that either the support has to be moved in two dimensions (WO 03/069391, DE 102 00 499), the detector has to be moved with respect to the support (US 2002/159057), the detector has to move in one dimension and the support in the other dimension or the optics has to include one or two dimensional scanning means i.e. galvo mirrors. On the other hand, the main difficulty of the strategy to illuminate the whole support simultaneously is to assure a homogeneous illumination across the whole distribution of sites. An alternative to the homogeneous illumination of the whole distribution of sites is the use of an array of light sources, whereby each site is illuminated by its own light source. DE 101 31 687 describes this strategy for the evaluation of PCR in a thermocycler with a plurality of wells using a beam splitter and an array of LEDs for illumination. DE 101 55 142 describes the dark field monitoring of fluorescence signals, wherein the microarray is illuminated by an array of LEDs, too, but in this embodiment no beam splitter is needed.\nConcerning the requirement to separate the optical path of the excitation beam and of the fluorescence light at least partially, there are again two different possibilities. The first possibility is the so called epi-illumination, whereby beam splitters are utilized and the excitation beam and the fluorescence light share at least part of the optical train. The second possibility is the use of oblique illumination. Here, the excitation beam is arranged in such a way that it has a certain angle to the normal of the support surface and the corresponding reflection of the excitation beam is outside of the acceptance angle of the detection system (e.g. US 2002/0005493 A1, EP 1 275 954 A2).\nUS 2003/0011772 A1 describes an optical apparatus to simultaneously observe a plurality of fluorescence dyes in a probe using a beam splitter. DE 197 48 211 A1 discloses a system to monitor the fluorescence signals generated in the wells of a microtiter plate simultaneously using a beam splitter, a field lens and an array of lenses focusing the light into each well. The detection is performed by imaging the light onto an array of photodiodes or a CCD chip. The fluorescence light collected in this embodiment of the system is appointed by the amount of dyes excited by the light cone of the focusing lens and therefore is dependent on the fill level of the well. WO 99/60381 claims an instrument for monitoring PCR reactions simultaneously in a plurality of vials in a temperature cycled block. The optical components of this instrument include again a beam splitter, a field lens, an array of vial lenses focusing individual light beams into each vial and a detection mean focusing the emission light onto e.g. a CCD detector. Due to the necessity of an array of vial lenses, the size and the lateral density of individual sites is limited. The JP 2002014044 describes a fluorometric apparatus to monitor fluorescence generated at a plurality of wells. The optical components comprise a beam splitter and a lens system to illuminate the wells collectively with light being parallel to the direction of the depth of the wells. However, the image forming optical system condenses the light onto a detection mean. U.S. Pat. No. 6,498,690 B1 discloses a method for imaging assays with an objective comprising a telecentric lens. U.S. Pat. No. 6,246,525 B1 claims an imaging device for imaging a sample carrier comprising a Fresnel lens.\nThus, it was the object of the present invention to provide an improved device for simultaneous monitoring of fluorescence signals from a lateral distribution of sites by optimizing the optical path towards homogeneous illumination and accurate detection. In one aspect of the present invention, the problem to be solved relates to improvements in monitoring multiplexed real-time PCR in a microtiter plate format."}
-{"text": "This invention relates to optically interconnecting opto-electric components on different integrated circuit (IC) chips, and also on the same IC chip, using a printed circuit board (PCB) on which the IC chip or chips are mounted.\nMany electronics systems, including computer motherboards, include one or more IC chips mounted on PCBs. A PCB provides a surface on which the IC chips are mounted, and also provides electrical interconnections between the IC chips.\nSignalling speed requirements between different IC chips in the same electronics system, and perhaps mounted on the same PCB, are ever increasing. In some cases, electrical signaling may not provide the needed, or desired, bandwidth, or may provide the bandwidth with costs. Some of the costs include a more complex design, in terms of multiplexing and demultiplexing the signals into multiple parallel lines. There may also be costs in terms of noise, both because the speed of the signaling may be nearing phyical limits and because of cross-talk between the parallel electrical interconnects.\nOptical signaling, as compared to electrical signaling, offers significantly higher bandwidth and eliminates, or greatly reduces, the noise problems inherent with electrical signaling. An example where bandwidth requirements are making optical signaling between IC chips in the same system increasingly attractive, and in fact may require optical signaling, is in computer motherboards. For example, signaling between processor IC chips and memory IC chips on the same motherboard are already in some systems two gigabytes per second, and will certainly only increase in the future.\nOptical signaling, however, poses design challenges not posed with electrical signaling. For example, optical signaling requires there to be an optical waveguide interconnection between the signal source and detector. In some cases, the optical interconnection between two IC chips within the same system has been provided with conventional optical fibers. However, this approach has its disadvantages. First, the optical fibers add cost to the system. Also, optical fiber connectors are typically large, and thus consume sometimes precious space, and the labor involved in providing connections to optical fibers is typically significant.\nBetter approaches to providing optical interconnects between IC chips within the same system, and even between opto-electronic components on the same IC chip, are therefore needed."}
-{"text": "The present invention relates to an emergency-stop circuit, which is an integral part of the typical industrial machine. More particularly, this invention relates to a centralized switching system and method for an emergency stop circuit.\nIn industrial equipment, the traditional emergency-stop circuit consists of a xe2x80x9cself-latchingxe2x80x9d relay that contains a number of closed (kill) switches which are connected in series, and when any one of the switches is opened, the relay is de-energized. Power is restored when all kill switches are closed, and a xe2x80x9cmotors-onxe2x80x9d momentary switch (e.g., push-button switch) manually closes the contacts of the relay. The relay contacts are the last link in the serial chain of switches that energizes the coil of the relay. It is self-latching in the sense that when the motors-on switch is released, the contacts are in the coil energizing circuit that keep them closed in the first place. The coil energizing circuit is referred to herein as the emergency-stop circuit.\nA robust, traditional circuit may have many kill switches in the emergency-stop circuit. These switches are typically distributed all over the machine. For example, lever-type switches are installed on door panels, so that power is killed (i.e., shut off) when one of the doors opens. This is referred to as the normally open configuration (NO), which means that the switch must be tripped to conduct. This kind of kill switch is the first to be defeated in practice. It is often taped or strapped closed so that a door may remain open during operation of the machine. (A common purpose for the defeat is debugging by a maintenance technician.) When there are several doors defeated in this manner located throughout a large machine, the probability is higher than desirable for a maintenance technician to inadvertently leave a switch defeated and return the machine to what will be unsafe use. Also, the cycle of taping/strapping and removal thereof causes wear and tear on the lever-type switch for which it was not designed.\nOther types of kill switches used in the industry include over-travel switches. These switches normally operate in the closed configuration (NC), which means that tripping of the switch opens the circuit. These switches include lever-type, magnetic, infrared, or the like. To defeat over-travel switches, the switches are temporarily removed, terminals jumpered, mounting screws loosened, and brackets are slid out of the way. This also creates opportunity for mistakenly leaving kill switches defeated (or misaligned) throughout the machine when it is returned to service.\nAnother example of a kill switch is an air pressure switch sensing an air line that delivers required air to an air bearing spindle. In a demonstrating test, or debug mode, the machine may be run without the spindle running (no air supplied or air temporarily unavailable). This requires the jumpering of the kill switch during such time. Afterwards, forgetting to re-enable the switch allows running of the spindle without air, which leads to hardware damage.\nEvidently, safe use of the traditional emergency-stop circuit requires experience and diligence on the part of the maintenance technician who attempts to temporarily bypass sections of the circuit in order to test or debug the system. Oversight due to distribution of the switches over numerous parts of the machine/device can cause him to forget to re-enable a kill switch before returning equipment back to duty.\nAdditionally, in order to test and debug, the technician must also disable certain devices whose power is controlled by the emergency-stop circuit. There is no straightforward, universal way to do this other than disconnecting the power to the device. This may be easy in some cases or not possible, very cumbersome, or unsafe in others.\nA final consideration for these testing and debugging methods is the time required for a technician to trace through a machine in order to determine where to disable a kill switch or where to disconnect power to a device. Additionally, managerial time may be spent generating documentation in order to aid the technician\"\"s task. This becomes apparent when one considers a factory floor that possesses a vast array of one-of-a-kind machines, all of which utilize some variant of the traditional emergency-stop circuit. Here, hypothetically, each circuit possesses essentially the same topology but utilizes different components that are located in different places and connected by a slightly different wiring scheme.\nIn spite of this, implementation of traditional emergency-stop circuits that are intrinsically xe2x80x9csafexe2x80x9d is certainly feasible and has been done for many years. There are reasons for the apparent success. It is a simple circuit, even though it is distributed throughout the machine. It well established. There are few components. But these are also the reasons why the circuit has not matured.\nTypically, experienced engineers are reluctant to add new parts and kill switches to the circuit in an effort to xe2x80x9ckeep it simple.xe2x80x9d In developing prototypes or one-of-a-kind machines, important kill switches such as a watchdog circuit and a computer ready are often omitted. Also, some kill switches having solid state outputs (e.g. NPN) do not fit into the serially connected topology. Each requires an extra part, such as an intermediate electro-mechanical relay, whose contacts are in the kill switch chain, and whose coil is controlled by the solid state output. Because of this, sensors employing solid state outputs are avoided, and their less reliable mechanical counterparts are used instead.\nEssentially, there is a mindset among skilled engineers concerning the altering of the traditional circuit\"\"s topology. Typically, the skilled engineer begins a new project assuming that he will use the traditional circuit. Valuable time is spent on other areas and is not devoted to re-engineering the architecture for the traditional circuit or evaluating its expanded role in the project. In fact, it is not obvious to the skilled engineer to change the traditional circuit in any way in order to add functionality that can be safely incorporated within it. Such functionality, if implemented, is therefore left to be distributed throughout the remainder of the system, intermingled with unsafe subsystems such as the computer.\nWhen implemented, for example, secondary outputs, such as amplifier xe2x80x9cenablexe2x80x9d or xe2x80x9cinhibitxe2x80x9d signals, are not usually incorporated into an emergency-stop circuit. If driven at all, a software program running on a computer having optically isolated digital outputs usually drives them. Furthermore, other feedback signals, such as xe2x80x9cstatusxe2x80x9d or xe2x80x9cfaultxe2x80x9d signals, are not used in emergency-stop circuits as kill inputs. This is generally because each signal is in a non-conducting state when the circuit is killed, which prevents the traditional circuit from restarting. If used at all, these feedback signals are likewise connected to the computer for the purposes of monitoring.\nDesigning in this way fosters subtle system-wide shortcomings, which can permit potentially unsafe or undesirable operation. Resulting failures or odd performance is not attributed to the emergency-stop circuit, since its simple circuitry and lack of substantial functionality are not directly responsible. Consequently, effort is typically not expended to evaluate its functionality.\nOne of the shortcomings becomes apparent when the traditional system enters into a power-loss period, which generally begins when the emergency-stop circuit is killed and ends when all residual power has been dissipated. During this brief period (e.g., 2 sec.), uncontrolled motion of motors can occur for some designs, because the motors are not being controlled, yet they are still technically powered by residual power in the system. In order to suppress this, designers have used the computer-controlled secondary outputs (enable, inhibit) in conjunction with the emergency-stop circuit to simultaneously cut power and disable the connected devices. This works in most cases, but is tedious to design, not flexible, and application specific. One case when this design fails is when the building power fails, which causes the computer to also cease functioning. Here the inhibit signal may not get to the device, which again creates an environment for briefly uncontrolled motion.\nMost of the examples found in existing technology are concerned with passive monitoring of the emergency-stop circuit. This approach is useful in determining which kill input was responsible for stopping the circuit, but it does not provide any configuration options for startup or power-loss periods. The following patents, each of which is incorporated herein by reference, demonstrate this approach: U.S. Pat. No. 4,263,647 to Merrell, et al, entitled xe2x80x9cFault Monitor for Numerical Control Systemxe2x80x9d; U.S. Pat. No. 5,451,879 to Moore, entitled xe2x80x9cElectromechanical Relay Monitoring System with Status Clockingxe2x80x9d; U.S. Pat. No. 4,616,216 to Meirow, et al., entitled xe2x80x9cEmergency Stop Monitorxe2x80x9d; and U.S. Pat. No. 5,263,570 to Stonemark, entitled xe2x80x9cConveyor Belt Emergency Stop Indicator Light System.xe2x80x9d Configuration options do exist in the above noted patents but only in the form of providing cascaded inputs and outputs so that multiple groups of sensors may be monitored. Other patents of interest include the following: U.S. Pat. No. 4,912,384 to Kinoshita, et al., entitled xe2x80x9cEmergency Stop Control Circuitxe2x80x9d discloses the traditional active portion of the emergency-stop circuit; U.S. Pat. No. 5,319,306 to Schuyler entitled xe2x80x9cPortable Electrical Line Tester Using Audible Tones to Indicate Voltagexe2x80x9d discloses circuits that provide audio status in the form of line testers, where the leads are brought into contact after the line is energized to check it.\nTraditional approaches to supplying power to motors during a power-loss period (period beginning with the loss of AC motor power and ending with either the total loss of all stored DC motor power or the loss of regulation of any associated logic power supply, whichever comes first) have focused on coarse (non-servo) control or decelerating motors to full stop. However, no approach exists that relates to fields employing emergency-stop circuitry.\nOther patents in this general field are also noted. For example, U.S. Pat. No. 5,278,454 to Strauss, et al. discloses an invention related to the heating, ventilation, and air conditioning field. It describes a motion control system that senses a loss of incoming power and utilizes a dedicated pre-charged circuit to act as a short duration power supply to effect gross motion of a motor to close a damper. U.S. Pat. No. 5,426,355 to Zweighaft, et al., entitled xe2x80x9cPower-Off Motor Deceleration Control Systemxe2x80x9d discloses an invention related to the tape drive industry in which a motion control system whose amplifier stores a dedicated internal PWM signal responsible for supplying open-loop deceleration commands for a given configuration of the tape drive system that is experiencing a power-loss period. U.S. Pat. No. 4,481,449 to Roda entitled xe2x80x9cPower Fail Servo Systemxe2x80x9d discloses an invention that also relates to the tape drive field which describes the use of several xe2x80x9cpower failxe2x80x9d signals that work in harmony to decelerate the motor towards full stop and uses the technique of dynamic braking to harness excess power in the storage capacitor. A signal exists in this example which monitors the logic power supply and appropriately disables (free wheels) the motor once the supply is out of regulation.\nThe present invention solves the problems in the art by providing a centralized programmable emergency-stop circuit that controls the flow of the power necessary for a machine to move its working elements. The invention possesses various levels of programmability that facilitate use of the same circuit across a wide variety of industrial applications and designs, as well as across a wide variety of operational scenarios for the same machine.\nThe circuit of the present invention includes various types of custom programmable kill inputs. These inputs are signals that, subject to their programming, can kill an energized emergency-stop circuit or prevent a killed circuit from energizing (startup). A given kill input can also be programmed to be ignored totally, to kill when inactive, or to also prevent startup when inactive. A given kill input can be programmed so that it only affects the energized circuit and does not restrict startup, and consequently, it may be inactive at startup. Such a programmed kill input is referred to herein as a xe2x80x9cfalling-type,xe2x80x9d because once it does go active, it is the active-to-inactive or falling transition that kills the circuit. Additional programming for the kill inputs exists such as digital filter parameters, clock selection, and the like, as well as time-out options for the falling-type kill inputs, which require them to go active within some period after startup.\nThe present invention also provides programming options to specify conditions for a motors-on signal to energize the circuit and for the control of secondary outputs. While the primary output of the circuit controls the flow of bulk power to working elements, it is the secondary outputs that connect in parallel to the working elements in order to inhibit or enable them. The method of programming secondary outputs determines their behavior, i.e., whether they are disabled entirely for the session, enabled only when the circuit is energized, or enabled based on one of the kill input signals. This latter setting permits a computer to keep a device enabled during a power-loss period, so that a reactionary movement can be effected which drains residual power left in the dying system.\nIn order to improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves performance during the period immediately following the application of electrical power needed to power circuit logic. Specifically, it is the object of the invention to inhibit energizing the circuit for a prescribed interval of time. Additionally, it is the object of the invention to provide programmability so that the interval may be changed.\nIn order to further improve performance during the period immediately following the application of electrical power needed to power circuit logic, it is the object of the invention to provide additional features and programmability. Specifically, it is the object of the invention to provide circuitry that determines whether the circuit has been energized at least once. Furthermore, it is the object of the invention to provide further additional circuitry that drives a dedicated power-up/reset error code which indicates electrical power has just been applied to the circuit logic. The power-up/reset error code therefore supersedes the conventional error code that is generated from all possible kill input sources. Additionally, it is the object of the invention to provide a clear signal capable of clearing the power-up/reset error code (so that the conventional error code may be revealed) and also capable of refreshing conventional error codes thereafter. It is also the object of the invention to provide programmability so that a set of clear input sources may be pre-selected from all available input sources.\nFinally, in order to further improve performance during the period immediately following the application of electrical power needed to power circuit logic, it is the object of the invention to provide additional features and programmability. Specifically, it is the object of the invention to employ a start signal that when inactive inhibits the initial energizing of the circuit. Activation of the start signal occurs in response to the final cycle of a specified number of deactivation and reactivation cycles of a ready-type input signal, and deactivation of the start signal occurs when the circuit is energized. Additionally, it is the object of the invention to provide programmability so that (1) the ability of the start signal to inhibit energizing is optional, (2) the specified number of cycles can be adjusted, and (3) a set of ready-type input signals may be pre-selected from all available input sources.\nIt is also the object of the invention to further employ the same start signal in subsequent energizing cycles in order to further improve performance. Specifically, a second specified number of deactivation and reactivation cycles is required in order to activate the start signal. Additionally, it is the object of the invention to provide programmability so that the second specified number of cycles can be adjusted.\nIn order to improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves how the circuit is commanded to energize. Specifically, it is the object of the invention to provide for additional nominal requirements for the activation of a motors-on signal, such as (1) requiring it to be previously inactive and (2) requiring it to be active for a prescribed interval or longer. Additionally, it is the object of the invention to provide programmability so that (1) the interval may be changed, (2) the requirement to be previously inactive is optional, and (3) a set of motors-on-type input sources may be pre-selected from all available input sources. Finally, it is the object of the invention to provide programmability so that (1) a set of monitor contact-type input sources may be pre-selected from all available input sources, where each monitor contact signal is active when the circuit is killed and the associated, downstream monitored relay has fully disengaged and (2) the requirement for a given monitor contact signal to be active for the motors-on signal to be active is optional.\nIn order to further improve the manner in which the circuit is energized, it is the object of the invention to employ a second start signal that when inactive inhibits the energizing of the circuit. Activation of the start signal occurs when all kill input sources are active, where programmability provides for a set of kill sources to be selected from all available input sources. Deactivation of the start signal occurs when the circuit is energized or when one or more of the kill input sources become inactive. Additionally, it is the object of the invention to provide status for the start signal. Furthermore, it is the object of the invention to accommodate watchdog-type kill input sources that toggle on-and-off repeatedly at a rate faster than a prescribed value, where the toggling is the requirement for the watchdog-type kill input to be active. It is also the object of the invention to provide programmability for this so that (1) the requirement for toggling is optional and (2) the minimum rate is programmable. Finally, it is the object of the invention to include in the generation of the start signal an additional, dedicated kill input source that indicates whether an internal circuit error exists.\nIn order to further improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves performance during the period immediately following energizing (right after it is started). Specifically, it is the object of the invention to provide audio status for a prescribed interval. Additionally, it is the object of the invention to provide programmability so that the interval may be changed.\nIn order to improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves the manner in which the circuit is de-energized (killed) or prevented from energizing. Specifically, it is the object of the invention to employ a kill signal that when active de-energizes the circuit or prevents it from energizing. Activation of the kill signal occurs when one or more kill sources become inactive, where programmability provides for a second set of kill sources to be selected from all available input sources. Deactivation of the kill signal occurs when all kill sources from the second set become active. Additionally, it is the object of the invention to include in the generation of the kill signal an additional, dedicated kill input source that indicates whether an internal circuit error exists.\nIn order to further improve performance for the manner in which the circuit is de-energized (killed) or prevented from energizing, it is the object of the invention to provide additional programmability so that pre-selected additional input sources can be dynamically added to the second set of kill sources at some point of time after the circuit becomes energized and subsequently removed at such time that the circuit is de-energized. A given, dynamically added input source may be programmed to be added immediately after the input source becomes active. Additionally, or alternatively, it can be added after a prescribed interval of time following the energizing of the circuit. It is also the object to provide programmability so that this prescribed interval can be adjusted.\nIn order to further improve performance for the manner in which the circuit is de-energized (killed) or prevented from energizing, it is the object of the invention to provided additional programmability so that one of the dynamically added input sources is dedicated to sensing the presence of the bulk power controlled by the circuit. Additionally, it is the object that this input source is an alternating-current type that generates a strobing signal indicative of the active state of the bulk power, where the strobing occurring at a rate faster than a prescribed value is the requirement that the kill input source is active. Finally, it is the object that the minimum rate is programmable.\nIn order to further improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves performance during the period immediately following de-energizing (right after it is killed). Specifically, it is the object of the invention to inhibit the re-energizing of the circuit for a prescribed interval of time after it is killed. Additionally, it is the object of the invention to provide programmability so that the interval for the dying period may be changed. Also, it is the object to provide audio or visual status during the dying period.\nIn order to further improve an emergency-stop circuit whose primary output controls the flow of bulk power needed for a machine to move its elements and whose secondary output controls the enable or inhibit of an element, it is the object of this invention to provide additional features and programmability for the circuit so that the source of the secondary output may be selected from a set of available sources. Specifically, it can be selected from the following sources: (1) none so that the element is always disabled, (2) from a signal that is active when the circuit is energized so that the element is enabled only when the circuit is energized, or (3) a dedicated enable-type input source, so that the element is enabled whenever the enable-type input source is active. It is also the object of the invention to provide additional programmability for the third case, which places a programmable pair of restrictions on when the enable-type input source has an effect so that it is used when (1) the circuit is energized or in the dying period that immediately follows de-energizing and otherwise, the element is disabled and (2) a watchdog-type input source is active and otherwise, the element is disabled. The requirement for the watchdog-type input source to be active is that it must toggle on-and-off repeatedly at a rate faster than a prescribed value. Finally, it is the object of the invention to provide additional programmability so that (1) the minimum rate for the watchdog-type input is programmable, (2) the enable-type input source may be pre-selected from all available input sources, and (3) the watchdog-type input source may be pre-selected from all available input sources.\nAccordingly, it is the object of the present invention to provide a programmable emergency-stop circuit that allows various options for the manner in which kill inputs affect the system and further provides options for the manner in which outputs are activated and deactivated. Furthermore, it is an object of the invention to provide programmability to specify the manner and timing for dynamically adding a given input source to the active set of kill inputs. Finally, it is an object of the invention to emp e circuitry that generally avoids software or a microprocessor, so that new functionality coupled with programmability may be safely incorporated within the emergency-stop circuit.\nOne important feature of the invention is its state machine, which provides a framework from which the invention operates. Defined by a set of internal signals that includes start and kill-type signals, the state machine specifies when the circuit may be energized, when it is killed, and when startup is inhibited. The internal signals are generated as a programmable function of time and input source states. Other features include audio status for startup and kill, requirements for startup that ensures desired energizing, requirements for a computer ready signal that ensures synchronization with software running on a computer, provisions for a dedicated error-code that identifies power glitches, and the safe oversight of a power-loss period during which a servo-controlled reflex action may be implemented.\nThe primary advantage for using the invention is that a centralized single circuit can be programmed and employed in a wide variety of machine designs. For a given machine design, for example, the circuit can be reprogrammed and thereby adapted to a different set of operational scenarios. When designing a machine or a plurality of machine/devices, the designer is able to associate any given input source with a desired kill input type that specifies how the input source affects the system. Furthermore, once operational in the field, for example, the machine will require maintenance, and to assist this, the circuit can be definitively reprogrammed from a central location so that certain inputs are temporarily but safely ignored and certain outputs are forced disabled during the maintenance operation.\nOther advantages of the invention are related to timing, filtering, and synchronization. One such advantage is the accuracy, and hence repeatability, that can be applied to timing the motors-on button\"\"s active period as well as to the timing of the start-up delay that prevents the immediate re-start during the DYING state of a freshly killed circuit. The use of timing and other related digital filters significantly reduces the susceptibility of the circuit to background noise. It is also an advantage from a system performance standpoint that the emergency-stop circuit causes the computer program and, thereby, the entire system to be in synchronization via several novel methods.\nThe invention will now be described, by way of example and not by way of limitation, with reference to the accompanying sheets of drawings and other objects, features and advantages of the invention will be apparent from this detailed disclosure and from the appended claims. All patents, patent applications, provisional applications, and publications referred to or cited herein, or from which a claim for benefit of priority has been made, are incorporated by reference in their entirety to the extent they are not inconsistent with the explicit teachings of this specification."}
-{"text": "1. Field of the Disclosure\nThe disclosure relates to a sense amplifier used in a semiconductor device, and particularly relates to a suitable sense amplifier in a semiconductor device that has a variable resistance memory cell, and to a data processing system.\n2. Description of Related Art\nConventional memory cells are known that store information based on the size of a resistance value or the \u201con\u201d current of a transistor. This type of memory cell generally has relatively high resistance values ranging from 10 k\u03a9 to several hundred kilohms (or kilo-ohm) even in a low memory state, and sense amplification is therefore usually performed using a highly sensitive differential current sense amplifier (see Japanese Patent Application Laid-Open No. 2004-39231)."}
-{"text": "1. Field of the Invention\nThis invention relates to a method of transferring data between computer systems with specific application in teleconferencing software programs. More particularly, it involves a method for transferring large amounts of data among interconnected computer systems according to the designated priority of the data, and for allowing the priority of the data to be changed before the data is completely transferred to a remote computer and for causing the remaining data to be subsequently transferred according to its new priority.\n2. Brief Description of the Prior Art\nWhenever two or more people are involved in the preparation of a document, whether it be a financial spread sheet, a CAD design, a circuit schematic layout, an organization report, a bit map image, etc., succeeding drafts of the document are prepared, circulated, and modified in the process. Each person annotates his or her remarks on the document and forwards it to the next person. Typically, several drafts of the document will be circulated before a final draft is produced. This is a very time consuming process.\nIn the case where a person involved in the document preparation process is at a different geographical location, getting the document from one location to another location and back becomes another tedious and time-consuming task. The document will either have to be mailed or faxed to that person, further complicating the entire process.\nOne standard method of alleviating this process is to hold meetings where everyone gathers and comments on the document with the hope of reducing the number of drafts needed before a final draft is produced. The shortcoming of this method is that there may be significant travel time and travel cost in getting all of the people to the same location. In addition, the final draft of the document is usually again circulated for final comments.\nOne solution to this problem is to use a teleconferencing software program, an aspect of which may contain an embodiment of the present invention. By using computer network connections or modem connected phone lines, everyone can be connected via his or her computer. By using the teleconferencing software program, everyone's computer screen displays the same data. In addition to using the software program and network or modem connections, conference calling over the voice phone lines or through the software program creates a dynamic and live atmosphere where everyone can participate in the discussion and refer to the document displayed on the screen.\nA very important capability of such teleconferencing software must be the ability to allow transfer of data from one computer user to other computer users. For example, in making a presentation using a number of frames of prepared graphs, charts, outlines, etc., each frame of data that is used must be quickly transferred to other users in the conference in order to have a common reference point for simultaneous discussion of the data presented. In addition, the presenter in the presentation may wish to skip among several frames of data or skip a few frames of data entirely. The teleconferencing software must allow this type of flexibility and still maintain a high efficiency in transferring data. At the same time, each frame of data must be organized in a manner that allows quick access by the users.\nAdditionally, the presenter may wish to transmit a private message to one particular user. The teleconferencing software will have to be able to distinguish between data for all interconnected computer systems (public data), and data for a particular user (private data), and properly transfer the data to the designated user or users.\nAnother problem in teleconferencing software is that the presenter may scroll through frames of data without allowing adequate time for the data to be transferred to all the computer systems. The presenter will eventually display one frame of data for discussion. At this time, this frame of data has the highest priority and must be immediately transferred to all other computer systems even though there may be several preceding frames of data that have not been completely transferred.\nNevertheless, all public frames of data scrolled through or loaded for the presentation must be transferred to all the connected users, because the presenter may eventually go back to previous frames of data in making his or her presentation. Thus, it is important to have the ability to organize the different frames of data and transfer the frames of data currently being used first, while establishing a system where other frames of data can be prioritized and transferred.\nAnother problem the present invention must deal with is the problem of transferring data between interconnected computer systems. The computer systems may be interconnected via modem, network, cellular links, or any other available connections. In connecting computer systems or nodes to computer systems, one computer may not be directly connected to all other computers involved in the teleconference. For example, referring to FIG. 1, there are four computer systems involved in this session of teleconferencing: computer A (10), computer B (12), computer C (14), and computer D (16). Computer A is only connected to computer B, computer B is connected to computer A and computer C, computer C is connected to computer B and computer D, and computer D is only connected to computer C. Computer A is connected to computer D only through computer B and computer C. In order for computer A to transfer data to computer D, the data must pass through computer B and computer C. Thus, if the user at computer A is making a presentation by using frames of data, these frames of data must travel through computer B and computer C to reach computer D. However, the data has to be transferred in such a manner so that there will not be a significant time lag between the time computer B receives the data, computer C receives the data, and computer D receives the data, so that all the users can follow the discussion or presentation in real-time or near real-time."}
-{"text": "1 . Field\nThe following description relates to a rectifier which may be used with a wireless power receiver.\n2 . Description of Related Art\nResonance power may include electromagnetic energy. A conventional resonance power transferring system may transmit power wirelessly, and may include a source device that transmits a resonance power and a target device that transmits a resonance power. Resonance power may be transferred from the source device to the target device.\nWhen an amount of current increases due to properties of a diode included in a conventional rectifier in a wireless power receiver (i.e., the target device of the wireless power transmission system), a voltage drop may increase due to resistance of the diode.\nVarious products, such as, for example, high-power applications that consume more than 100 W power and low-power applications that consume less than 10 W, have been studied. However, it has been found that for a wireless power transmission system that consumes about 10 W, the total efficiency is low, for instance, only about 60%."}
-{"text": "This section is intended to introduce the reader to various aspects of art that may be related to various aspects of the present disclosure, which are described and/or claimed below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present disclosure. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nA vehicle that uses one or more battery systems for providing all or a portion of the motive power for the vehicle can be referred to as an xEV, where the term \u201cxEV\u201d is defined herein to include all of the following vehicles, or any variations or combinations thereof, that use electric power for all or a portion of their vehicular motive force. As will be appreciated by those skilled in the art, hybrid electric vehicles (HEVs) combine an internal combustion engine propulsion system and a battery-powered electric propulsion system, such as 48 volt or 130 volt systems. The term HEV may include any variation of a hybrid electric vehicle. For example, full hybrid systems (FHEVs) may provide motive and other electrical power to the vehicle using one or more electric motors, using only an internal combustion engine, or using both. In contrast, mild hybrid systems (MHEVs) disable the internal combustion engine when the vehicle is idling and utilize a battery system to continue powering the air conditioning unit, radio, or other electronics, as well as to restart the engine when propulsion is desired. The mild hybrid system may also apply some level of power assist, during acceleration for example, to supplement the internal combustion engine. Mild hybrids are typically 96V to 130V and recover braking energy through a belt or crank integrated starter generator. Further, a micro-hybrid electric vehicle (mHEV) also uses a \u201cStop-Start\u201d system similar to the mild hybrids, but the micro-hybrid systems of a mHEV may or may not supply power assist to the internal combustion engine and operates at a voltage below 60V. For the purposes of the present discussion, it should be noted that mHEVs typically do not use electric power provided directly to the crankshaft or transmission for any portion of the motive force of the vehicle, but an mHEV may still be considered as an xEV since it does use electric power to supplement a vehicle's power needs when the vehicle is idling with internal combustion engine disabled and recovers braking energy through an integrated starter generator. In addition, a plug-in electric vehicle (PEV) is any vehicle that can be charged from an external source of electricity, such as wall sockets, and the energy stored in the rechargeable battery packs drives or contributes to drive the wheels. PEVs are a subcategory of electric vehicles that include all-electric or battery electric vehicles (BEVs), plug-in hybrid electric vehicles (PHEVs), and electric vehicle conversions of hybrid electric vehicles and conventional internal combustion engine vehicles.\nMicro Hybrid technology can use a dual voltage architecture, such as a traditional 12V vehicular electrical system used in conjunction with a lead-acid battery, and a 48 volt vehicular electrical system used in conjunction with a Lithium-ion battery. 12 volt electrical system, as used herein, refers to a traditional vehicular electrical system that operates at a nominal 12 volts. The actual voltage varies dynamically depending in part on the charge state of the battery and the load, and an any point in time can be more or less than 12 volts. 48 volt electrical system, as used herein, refers to a vehicular electrical system that operates at a nominal 48 volts, such as one using an LI-ion battery. The actual voltage varies dynamically depending in part on the charge state of the battery and the load, and an any point in time can be more or less than 48 volts. The 12 volt system can include things such as lights, audio/entertainment, electronic modules and ignition. The 48 volts system can include the A/C compressor, active chassis, and regeneration. These systems support higher power loads and provide redundancy. Typically an 8-10 kW motor/generator captures energy for regeneration, supports re-start and supports higher power loads. A DC/DC converter bridges between the higher 48 volt system and the traditional 12 volt system.\nSuch a micro hybrid vehicle can change electrical load management due to high power regeneration, and provide for electrification of new loads such as air conditioning, active chassis and safety, electric supercharging, as well as result in increased fuel efficiency.\nThe DC-DC converter needed for to bridge the systems should be able to provide sufficient power without taking excess space. Moreover, it should be able to withstand the vehicular environment, including high temperatures."}
-{"text": "Migraine is a chronic condition with recurrent episodic attacks. It is rather unpredictable illness with its characteristics varying among patients. This unpredictability and variability is also observed within migraine attacks observed in a single patient. Among the most distinguishing features of a migraine is a potential disability caused by the accompanying headache and nausea with or without vomiting as well as extreme sensitivity to sound and light (Headache, 39: 720-727 (1999)). Because of the variability and complexity of the condition, effective management of patients suffering from migraines is challenging.\nMigraine headaches which are considered \u201cprimary headaches\u201d are about three times more common in women than in men. Geographically, the occurrence of migraine headaches varies significantly and ranges from 1.5% in Southeast Asia to 14% in Western countries (GRIM, Cephalalgia, 12:229 (1992); JAMA, 267:64-69 (1992) and Pharmacoeconomics, 11: 1-10 (Suppl.1)(1997)).\nSystemic administration of anti-migraine and anti-nausea drugs orally to patients has not been very successful, in part because migraine is often accompanied by nausea and the orally administered drugs are vomited before they can take effect. The only viable route of administration for treatment of nausea and/or migraine is the intravenous or another injectable administration. These typically require a visit at the doctor's office or hospital. The failure to successfully treat migraine or nausea is thus based on a delivery method rather than on the drug effectiveness.\nThe vaginal delivery route of drugs through the vaginal mucosa to the uterus and/or to the general circulation has been discovered by inventors and is disclosed, for example, in the U.S. Pat. Nos. 6,086,909, 6,197,327 and 6,572,874 and in a co-pending application Ser. No. 10/600,849 and 10/349,029, all hereby incorporated by reference.\nAs well as the vaginal delivery route described in the above cited patents and application works, there is still some need for improvement, particularly as it concerns an efficacious quantitative drug delivery.\nThe current invention thus concerns an improved transmucosal delivery of anti-migraine and anti-nausea drugs through vaginal mucosa directly to uterus or to the general circulation which is more efficacious due to a more quantifiable sequestration of the drug within the impermeable layer or layers covering a proximal portion of the vaginal device.\nIt is therefore a primary objective of this invention to provide a vaginal device, such as a tampon, tampon-like foam or another vaginal device which is coated, or of which a proximal portion is coated, with a fluid impermeable layer(s) of film, foil, foam or xerogel forming a strip or an attached or removable cap or cup wherein said impermeable layer further comprises a mucoadhesive composition comprising an anti-migraine or anti-nausea drug, or a combination of both, said composition being released and delivered from said layer(s) substantially quantitatively through the vaginal mucosa into the uterus and/or to the general circulation."}
-{"text": "I. Field of the Invention\nThis invention relates generally to control of machines and particularly to eliminating relative misalignment of a plurality of driving means spaced along the length of a moveable member, the member being propelled by the driving means to move linearly transverse to its length and the misalignment being measured parallel to the direction of motion.\nII. Description of Related Art\nAn example of machine construction giving rise to misalignment of driving means of a machine member, referred to as skew, is illustrated in FIG. 1. Machine member 20 is moveable along rails 24 and 26 as indicated by the double ended arrow \"X\". Skew is defined as a difference of the positions of the driving means or driven points of machine member 20 as measured relative to a common reference and parallel to the axis of motion. For example, in FIG. 1 skew is represented by a difference in the values of X' and X\". Movement of member 20 is imparted by, respectively, actuators 28 and 30 driving through transmissions or gear trains 32 and 34, respectively, at distal ends 21 and 22. It will be appreciated that machine constructions may include additional driving means spaced along the length of member 20. The driving means may include, for example, pinions engaging racks on the rails to propel member 20. Other driving devices may be included in the driving means for driving member 20 such as, for example, nuts rotated by actuators and engaging screws mounted parallel to the rails.\nMember 20 and driving means 28 and 30 are arranged to propel member 20 linearly along the rails 24 and 26, the rails being substantially parallel to one another and substantially perpendicular to member 20. Skew may introduce binding of member 20 against rails 24 and 26 as well as poor engagement of the driving means. Skew can result in unsatisfactory response of actuators 24 and 26 to the control thereof and may result in excessive mechanical wear or damage to member 20, rails 24 and 26 or the driving means. Skew may arise as a result of forces acting on member 20 during periods when actuators 28 and 30 are not energized and therefore not active to maintain positions of ends 21 and 22. In light of the adverse consequences of skew, it is desirable to eliminate skew before initiation of extensive motion of member 20.\nIt is known from U.S. Pat. No. 4,045,660 to return a machine member to a desired location following power interruption using values of measured position determined upon power loss and upon power restoration. It is known from U.S. Pat. No. 4,484,287 to restore a moveable machine member to a desired position following power interruption by storing position information in a nonvolatile memory upon power loss for recall upon restoration of power. From U.S. Pat. No. 5,013,988 it is known to use presettable counters in conjunction with absolute encoders to create and update absolute position data, the data being updated so long as power is applied to the measuring system. From U.S. Pat. No. 4,629,955 it is known for control of motion of a machine member driven at distal ends to vary servomechanism gain of the driving means to eliminate skewing therebetween. The control of this patent provides only for reduction of skew attributable to differences in loading on the driving means during execution of motion. Servomechanism position control effected in accordance with this patent does not provide for detection, reduction or elimination of skew between the driving means upon application or restoration of power.\nThe controls of the aforesaid references addressing restoration of position after power interruption all rely on absolute position information which is immediately available upon application or restoration of power. However, in the event sufficient absolute position data is not available on application of power, it must be determined from position transducers. Sufficient position data will not exist if, during a period when position control is disabled, positional changes are not monitored and the range of position measurement transducers is less than the potential magnitude of positional change or the output of the position transducers while the machine member is stationary do not provide any indication of position whatsoever. Position transducers of the latter type include incremental or semi-absolute encoders, which are favored for providing high resolution position measurement over extended distances."}
-{"text": "The present invention relates to a braid and a braiding machine, and more particularly, to a polygonal braid and a braiding machine therefor, in which a guide plate having a track capable of braiding a polygonal braid and a feed gear corresponding thereto are provided. The polygonal braid braided in a square or triangular shape can be utilized as a binding string for shoes or clothes etc because of increased binding force by a polygonal. Braids of various quality and colors may be created by using different qualities or colors of strands.\nBraids are utilized in several fields, for example, as part of an electric wire or hose, as a binding string etc. A specific braid is formed on the outer circumference of the electric wire or the string and provides an elastic and relaxed covering for an interior electrc wire or a string etc. and protects the interior electric wire from being contaminated or damaged by impact, braids are often used in place of string for daily use in shoes or clothes etc., in addition to specialized uses.\nA general braider is composed of a guide plate having a track on which a spool is moved, a feed gear for moving an electric spool along the track the guide plate, a driving gear for driving a plurality of feed gears and a plural number of rollers on which a braided wire is wound, etc.\nFIG. 6a shows a guide plate for manufacturing a general cylindric braid and its braid. On this guide plate 100, two tracks 101, 101xe2x80x2 on a gentle circular curve line of a jig jag shape are formed, intersected with each other. As shown in FIG. 6b, on its lower part, a plurality of feed gears 102, 102xe2x80x2 are positioned beneath and aligned within the intersected curves of tracks 101, 101xe2x80x2. In such construction, the plurality of feed gears 102, 102xe2x80x2 are simultaneously rotated by a rotation of the driving gear 103, yarn from separate spools are combined within feed gears 102, 102xe2x80x2 onto the guid plate 100 which is rotated.\nTherefore, a plural number of spools of yarn are rotated, repeatedly performing a rotation and a revolution centering around a center point of the guide plate 100, feeding out yarn, which are intersected with one another, rotating along the track 101, 101xe2x80x2.\nOn an outer circumference of the central yarn, thus, a braid based on a cylindric shape is produced by the rotation of the spool as shown in FIG. 6c. \nThe ordinary cylindrical braid as described above, when used as a binding of shoes, has a low frictional force due to the cylindrical shape which can result in the shoe lace coming loose.\nWhan a braid is made using a single color, by prior art methods, the brain color is monotonous; when using a single color, by prior art methods, the braid color is monotonous; when using various colored braid, the braid color may appear untidy.\nIn order to overcome the problems presented in prior art braiding methods, the present invention teaches a braid formed in a polygonal shape such as a triangle or a square shape with each edge of each face the polygon intersecting with an edge of the adjoining face of the polygon."}
-{"text": "Some conventional vehicles include an electrical wiring harness extending between a roof and a roof liner of the vehicle. The electrical wiring harness is attached to the roof liner using tape and also using two-piece clips. A first piece of each clip is glued to the roof liner and a second piece of each clip is taped to the electrical wiring harness and is snapped into the respective first piece."}
-{"text": "It is well recognized in the petroleum industry that boron containing compounds are desirable additives for lubricating oils. One such boron containing compound is disclosed in U.S. Pat. No. 3,224,971 to Knowles et al. which relates to intracomplexed borate esters and to lubricating compositions containing said esters. The borate esters are organo-boron compounds derived from boric acid and a bis (o-hydroxy-alkylphenyl) amine or sulfide.\nAnother extreme pressure lubrication composition is disclosed in U.S. Pat. No. 3,185,644 to Knowles et al., which relates to lubricating compositions containing amine salts of boron-containing compounds. The amine salts are formed by reaction of a hydroxy substituted amine and a trihydrocarbyl borate. The amine-borate compounds thus formed are described as useful as load carrying additives for mineral and synthetic base lubricating oils.\nBoric-acid-alkylolamine reaction products and lubricating oils containing the same are disclosed in U.S. Pat. No. 3,227,739 to Versteeg. These amine type products are prepared by reacting equal molar proportions of diethanolamine or dipropanolamine and a long chain, 1, 2-epoxide. The intermediate reaction product thus produced is reacted with boric acid to produce the final reaction product. These compounds are added to lubricants to prevent rust formation.\nAnother boron ester composition is described in U.S. Pat. No. 3,269,853 to English et al. which discloses a boron ester curing agent which consists of a cyclic ring structure containing boron, oxygen, nitrogen, carbon and hydrogen.\nAnother boron composition is disclosed in U.S. Pat. No. 3,598,855 to Cyba which relates to cyclic borates of polymeric alkanolamines formed by reacting a borylating agent with a polymeric alkanolamine. The compounds thus formed are described as additives for a wide variety of petroleum products including lubricating oils.\nCurrently, there are phosphorus-containing additives which provide extreme pressure, anti-wear and/or friction-reducing properties to automotive engine oils. However, with the advent of the catalytic converter, alternative additives are needed. During combustion in an automotive engine, any oil which leaks or seeps into the combustion chamber yields phosphorus deposits which poison the catalyst in the catalytic converter. As a result, there is a need for automotive engine oil additives which are phosphorus-free but provide useful extreme pressure, anti-wear, and/or friction-reducing properties to the oil.\nAccordingly, it is one object of the invention to provide a phosphorus-free additive having such properties and which, upon combustion, will not adversely affect the catalyst in the automotive catalytic converter.\nIt is yet another object of the present invention to provide boron-containing, heterocyclic compounds or derivatives thereof which have extreme pressure, anti-wear and friction-reducing properties.\nYet another object of the present invention is to provide a lubricating composition having extreme pressure, anti-wear and friction-reducing properties.\nA further object of the present invention is to provide a lubricating composition containing extreme pressure, anti-wear, friction-reducing and corrosion prevention additives, and in addition, an anti-oxidant to prevent attack of oxidants upon metal bearings.\nOther objects and advantages of the invention will be apparent from the following description."}
-{"text": "Developmental disorders such as autism spectrum disorders (ASD) affect nearly 14% of children in the United States. Diagnostic methods for conditions such as ASD vary considerably, and even the use of \u201cbest practice\u201d tools provides rather poor sensitivity and specificity to the conditions. Late diagnosis of developmental disabilities reduces effectiveness of treatments and often results in poor outcomes. Furthermore, treatment providers (e.g., pediatricians or other medical professionals) lack adequate tools for measuring progress in these conditions."}
-{"text": "Makeup cases of this type generally contain a grille over a shallow cup of makeup product of powder cake, as well as an applicator puff enclosed between the grille and the cover. These cases can only be used dry.\nThere are also known cases of the same type containing a moistenable sponge for moist use."}
-{"text": "The present invention relates to a monomeric composition and a polymer obtained by the polymerization thereof or, more particularly, to a monomeric composition capable of being polymerized into a curable polymer which gives a rubbery elastomer having excellent heat and cold resistance and oil resistance as well as a polymer obtained from the monomeric composition.\nSo-called acrylic rubbers belong to a class of synthetic rubbers obtained by the copolymerization of an acrylic monomer such as ethyl acrylate, n-butyl acrylate, 2-methoxyethyl acrylate, 2-ethoxyethyl acrylate and the like with a comonomer which gives crosslinking points in the molecules of the copolymer. Acrylic rubbers obtained from ethyl acrylate as the principal comonomer have excellent oil resistance and heat resistance but the cold resistance thereof is poor. Acrylic rubbers obtained from n-butyl acrylate as the principal comonomer have excellent heat and cold resistance but the oil resistance thereof is poor. Further, acrylic rubbers obtained from 2-methoxyethyl or 2-ethoxyethyl acrylate as the principal comonomer have excellent cold resistance and oil resistance but they have rather poor heat resistance. Thus, none of the conventional acrylic rubbers satisfies the requirements for the oil resistance and cold resistance simultaneously."}
-{"text": "Firearm marksmen, particularly military sharp shooters, have a need for supporting the forward end of a firearm in a stable adjustable manner. Often, a bipod support is used for such front end firearm support. Military sharp shooters have a particular need for a portable, light weight and retractable bipod which also offers significant degrees of adjustability. In particular, it would be useful to have a bipod support having pivotably mounted legs wherein the legs may be adjusted to various positions including a retracted position in which the legs are generally parallel to the longitudinal axis of the firearm. It would also be useful for the legs of such a bipod to have adjustable telescoping portions for adjusting the length of the legs. Moreover, it would be useful if such a bipod support were adapted to allow pivoting adjustment about a vertical axis and a horizontal axis with respect to the legs of the bipod for aiming adjustment."}
-{"text": "This invention relates to an improved elevator speed control apparatus for regulating the running speed of an elevator cage to accommodate load changes when passengers exit before the cage comes to a complete halt at an accessed floor.\nA conventional elevator speed control system is shown in FIG. 1, wherein an electric power converter 2 which comprises a plurality of thyristors connected in a 3-phase bridge configuration, is coupled to a 3-phase AC power source 1 and generates DC power that is supplied to an armature 3 of a DC elevator drive motor through a line 2a. The field winding for the motor is not shown in the drawing.\nA tachometer generator 4 driven by the armature 3 produces a speed signal on line 4a proportional to the rotation speed of the armature. A traction sheave 5 also driven by the armature 3 drives an elevator cage 7 and a counterweight 8 through a main cable 6 as is well known. A speed arithmetic circuit 10 receives the speed signal on line 4a from the tachometer generator 4 and a speed instruction signal on line 9a from a speed instruction signal generator 9 as inputs, and generates a current instruction signal on output line 10a. The speed arithmetic circuit 10 along with the speed instruction signal generator 9 and the tachometer 4 constitute a speed control system.\nA current arithmetic circuit 12 receives as inputs the current instruction signal on line 10a from the speed arithmetic circuit 10 and a current signal on line 11a from a current detector 11 proportional to the current supplied to the converter 2. A phase shifter 13 receives the output signal on line 12a from the current arithmetic circuit 12 as an input, and outputs a firing control signal on line 13a for the converter 2. The current arithmetic circuit 12 along with the current detector 11 constitute a current control system.\nBy controlling the firing angle or phase of the thyristor converter 2 by means of both the speed control and current control systems, the voltage applied to the armature 3 is correspondingly controlled and thus the running speed of the elevator cage 7 is controlled through the traction sheave 5. In other words, the elevator cage 7 is speed controlled in accordance with the difference between the speed instruction signal on line 9a and the actual speed signal on line 4a with a high degree of accuracy.\nIn the aforementioned speed control system, in order to compensate for the non-linearity of the converter 2, the response time of the minor loop constituted by the current control system is set at an extremely short value, generally in the range of 0.01 to 0.03 second. On the other hand, the response time of the main loop constituted by the overall speed control system must be set at a higher value in order to avoid resonances in the suspension and traction cables. Consequently, the speed control system is generally designed so as to have a response time in the range of 0.2 to 0.33 second.\nWith such a conventional elevator system, in order to improve the transport efficiency and speed up the overall operation both the internal cage door and the external door on the accessed floor are sometimes controlled to be simultaneously opened just before the cage reaches the floor. A brake system (not shown) is also provided to engage the traction sheave 5, but such engagement does not occur until the cage comes to a complete stop. Passengers may thus step out of the cage before the brake system acts upon the traction sheave, and as a result an abrupt variation in torque is exerted on the sheave due to the change in the cage load or weight, as shown in FIG. 2(a).\nUpon the occurrence of a torque variation, the current flowing through the armature 3 correspondingly varies in response to the output of the speed arithmetic circuit 10 due to the functioning of the current control system, as described above. In this case, however, based on the relatively slow response time characteristics of the speed control system the current flowing through the armature 3 varies or adjusts relatively slowly as shown in FIG. 2 (b). The running speed of the cage 7 therefore varies as shown in FIG. 2 (c), as a result of which the cage may overshoot or undershoot the exact position of the accessed floor, which constitutes a potentially dangerous situation. Even in the best case where the cage ultimately stops at the exact position of the floor sill, the passengers will experience a discomforting \"acceleration-deceleration bump\".\nIt will be understood that the curves of FIG. 2 have been simplified by removing or subtracting therefrom the normal transient values, to leave just the \"abnormal\"variants caused by a premature passenger exit (or entry)."}
-{"text": "The present, invention relates to a manufacturing method of a combination material of metal foil and ceramic by joining a metal foil onto surfaces of various types of ceramics, and also relates to a metal foil laminated ceramic substrate manufactured from said combination material of metal foil and ceramic.\nRecently, a new improvement has been attempted by combining ceramics with metals in the application fields of ceramics. For example, a metal material composed of nickel alloy, titanium alloy, chromium alloy, or the like is joined to a ceramic material composed of alumina, zirconia, magnesia, or the like by using a combination technique such as diffusion combination so as to manufacture a combination material of metal and ceramic, which is used in various devices and apparatuses.\nHowever, since such conventional combination material of metal and ceramic is manufactured by simply joining a metal material to a ceramic material under a certain pressure, when the pressure is applied, there sometimes occurs fracture of the ceramic material which is a much more brittle material as compared with the metal material, thus causing problems in productivity.\nAn object of the present invention is to solve the problem mentioned above and provide a manufacturing method of a combination material of metal foil and ceramic capable of completely joining a metal foil to a ceramic material even with a low pressure applied thereto and preventing the ceramic material from fracturing so as to improve the productivity of a combination material of metal foil and ceramic or a metal foil laminated ceramic substrate, and also provide such metal foil laminated ceramic substrate.\nIn order to achieve the object mentioned above, the manufacturing method of a combination material of metal foil and ceramic comprises the steps of ion-etching a surface of a metal foil and a ceramic material to be joined together to activate and clean the surfaces, heating said surface of the ceramic material, which is held on a holder, to a temperature range of 250 to 500xc2x0 C., pressure-welding said surface of the metal foil to said surface of the ceramic material held on the holder under a pressure not more than 1 kg/mm2, and thus heat-joining the metal foil and the ceramic material to manufacture a combination material of metal foil and ceramic."}
-{"text": "Biomass thermal conversion is an attractive method for generating synthetic gas to run engines or to produce useful end products such as charcoal. Carbonaceous byproducts are typically inexpensive or free to source. Unfortunately, biomass byproducts come in a wide array of shapes and sizes, and extra machinery is usually required to pre-process the feedstock into forms acceptable to gasification or pyrolysis machines. This processing equipment is often expensive and difficult to operate, which challenges the ultimate attractiveness of biomass thermal conversion solutions.\nThus, there is a need in the field of biomass thermal conversion for system capable of utilizing a wide range of fuel shapes and sizes, without feedstock preprocessing on the front end. This invention provides such a solution through a novel \u201creactor-internal\u201d fuel processing solution that reduces a wide range of input biomass feedstock to a common size of granulated char."}
-{"text": "1. Technical Field\nThe present invention relates to data processing and, in particular, to scripts in a network data processing system. Still more particularly, the present invention provides a method, apparatus, and program for evaluating scripts.\n2. Description of Related Art\nThe worldwide network of computers commonly known as the \u201cInternet\u201d has seen explosive growth in the last several years. Mainly, this growth has been fueled by the introduction and widespread use of so-called \u201cweb browsers,\u201d which enable simple graphical user interface-based access to network servers, which support documents formatted as so-called \u201cweb pages.\u201d These web pages are versatile and customized by authors. For example, web pages may mix text and graphic images. A web page also may include fonts of varying sizes.\nA browser is a program that is executed on a graphical user interface (GUI). The browser allows a user to seamlessly load documents from the Internet and display them by means of the GUI. These documents are commonly formatted using markup language protocols, such as hypertext markup language (HTML). Portions of text and images within a document are delimited by indicators, which affect the format for display. In HTML documents, the indicators are referred to as tags. Tags may include links, also referred to as \u201chyperlinks,\u201d to other pages. The browser gives some means of viewing the contents of web pages (or nodes) and of navigating from one web page to another in response to selection of the links.\nBrowsers may also read and interpret pages including scripts, such as JAVAScript or JScript. JAVAScript is a popular scripting language that is widely supported in Web browsers and other Web tools. JAVAScript adds interactive functions to HTML pages, which are otherwise static, since HTML is a display language, not a programming language. JScript is similar to JAVAScript, but has extensions specifically for the Microsoft Windows environment.\nDifferent Web browser software applications support scripts to different degrees. In fact, different versions of a browser application may support scripts differently. The rivalry between browsers in the market, such as Netscape Navigator by Netscape Communications and Internet Explorer by Microsoft Corporation, has led to a disparity between standards. Netscape uses different syntax with JAVAScript than Internet Explorer uses with JScript.\nBrowser and platform dependent script functions can cause a nightmare for script development, testing, and maintenance. Therefore, it would be advantageous to provide an improved script evaluator for determining browser support."}
-{"text": "1. Field of the Invention\nThis invention is related to laser systems. In particular, this invention deals with adjusting alignment of laser beams in laser systems.\n2. Description of the Related Art\nWhen fabricating memory circuits, a laser repair system can be used to selectively sever conductive links, effectively removing faulty memory cells from the circuit.\nAs the size and spacing of link elements decreases, laser repair systems have had to increase in accuracy in order to perform their intended function. The complexity of a laser repair system capable of such accurate operation is significant. Multiple mirrors and other optical elements are used to generate and position a laser beam spot for severing a conductive link. Like the circuit fabrication process itself, laser repair systems are subject to many complex factors. For example, thermal expansion may lead to changes in the orientation or position of optical elements in the path of a laser beam. These changes to the elements that affect the laser beam can cause the laser beam spot to drift away from its intended location and can cause errors when trying to repair a circuit. Although the beam spot position is aligned with reference to wafer alignment markers with every new wafer processed, a misaligned laser beam path that deviates from a normal orientation to the work surface can still produce beam spots of unintended location, shape and/or size which adversely affect operation of the repair system.\nU.S. Pat. No. 6,483,071 (hereinafter referred to as the '071 patent) entitled \u201cMethod and system for precisely positioning a waist of a material-processing laser beam to process microstructures within a laser-processing site\u201d is assigned the assignee of the present invention. The disclosure of the '071 patent is hereby expressly incorporated by reference in its entirety. The '071 patent discloses many features of a laser based system for memory repair, and is particularly related to accurate (sub-micron) and high-speed positioning of a laser beam waist relative to a link or similar target structure. In the '071 patent, an air-bearing based assembly was disclosed for positioning of optical components (e.g: an objective lens) along the optical (Z) axis. In addition to noise and reliability issues (ie: wearing mechanical parts) it was recognized that X,Y displacement errors during Z axis motion are much better controlled or eliminated with an air bearing system. Such displacements, even if a fraction of a micron, can lead to link severing results which are incomplete (e.g. contamination) or possibly cause damage to surrounding structures. Hence, a displacement of a laser beam from a target location by a fraction of one-micron, corresponding to a fraction of one spot diameter, may generally lead to reduced yield.\nTraditionally, laser repair systems have undergone periodic, manual adjustment to correct problems with alignment. For example, every month, a trained technician may have to manually adjust optical elements in order to correct alignment problems that have developed since the last adjustment. In the M430 laser link blowing machine from GSI, coarse adjustments to laser beam alignment were made by manually adjusting the laser beam orientation while viewing the laser beam spot with a \u201cthru-lens viewing system\u201d (TTLV). The TTLV is essentially a camera and TV monitor arrangement coupled to the laser beam path. The spot position was determined relative to a crosshair. The beam was first aligned to be centered in the lens aperture. Then the beam was aligned for zero spot translation during zoom expansion. Zoom adjustments corresponded to a range of spot sizes. If the beam was properly aligned along the Z-axis, the beam would appear stationary on the monitor for all zoom settings. Finer beam alignment was carried out by adjusting the spot size to a minimum, placing a calibration grid on the work surface, and performing iterative manual adjustments of turning mirrors to align the optical system and reduce any lateral (X-Y) displacement to within a specified tolerance.\nThis traditional approach to adjusting the alignment of a laser beam has several drawbacks. For example, the means used by the technician to determine beam alignment may itself be subject to error. Alignments based on erroneous alignment data may augment alignment problems in the system. Other problems may include the significant time expense involved in manual adjustment. Delays arising from manual alignment can represent a serious cost for businesses operating laser repair systems. For these reasons and others, automated methods of static laser beam alignment have been developed. Such methods are described for example in U.S. Pat. Nos. 5,011,282 to Ream, et al., 5,315,111 to Burns, et al., 5,923,418 to Clark et al., and 6,448,999 to Utterback et al. Of these prior patents, Burns, Clark, and Utterback split off portions of the laser beam to optical detectors placed adjacent to the laser beam path. Alignment of the beam with respect to the detectors is used to deduce alignment of the beam to the workpiece. In the '282 patent to Ream, changes in laser beam spot position on a target are used to determine a laser beam deviation angle, which can then be used to correct the laser beam path alignment."}
-{"text": "Modern application-specific integrated circuits (ASICs) integrate greater and greater security and data protection functionality into the hardware (HW). The integrated functionality provides more reliable and more efficient hardware security for both conventional \u201cData At-Rest\u201d and conventional \u201cData In-Flight\u201d protection.\nData storage systems are moving to distributed storage models that are based on storage networking. The move has an impact for enterprise data protection: the distributed models increase the vulnerability of stored data (i.e., Data At-Rest) to various attacks, both external and internal and both malicious and accidental. For Internet traffic and other moving data (i.e., Data In-Flight), the move provides such protection as sender and recipient mutual authentication, key exchange, data confidentiality, authenticated encryption (which is a type of encryption/decryption that additionally providing a way to check data integrity and authenticity) and replay protection.\nIn contemporary applications, the speed/throughput of the traversing data is up to 10 Gb/s (gigabits per second) and beyond. For some storage applications, the speed/throughput of the traversing data is even 10\u00d7 higher: up to 100 Gb/s and beyond. The high speeds alone make security support of the data in software (SW) almost infeasible as far as security transformations are usually incorporated into the main data path and appear as bottlenecks from efficiency and performance standpoints.\nMany cryptographic protocols use an encryption process and message authentication and data integrity services independently with each process using an independent key. To speed up overall computations, new cryptographic modes that combine and provide both crypto services using a single \u201ccombined\u201d mode were proposed and became accepted by both the National Institute of Standards and Technology (NIST) and the Institute of Electrical and Electronics Engineering (IEEE) and other technical professional organizations and committees working in network and data storage security areas.\nTo prevent data lost and breach, IEEE P1619 \u201cStandard Architecture for Encrypted Shared Storage Media\u201d suggests using the XTS-AES (Advanced Encryption Standard) (XOR-Encrypt-XOR (XEX)-based Tweaked Electronic Code Book (ECB) mode with Cipher Text Stealing (CTS)). The P1619.1 \u201cStandard for Authenticated Encryption with Length Expansion for Storage Devices\u201d uses the Galois/Counter mode (GCM), Counter mode (CTR) with Cipher-Block Chaining (CBC)-Message Authentication Code (MAC) (CCM) and other cryptographic processes. Both drafts are now accepted standards: IEEE Std. 1619-2007 and IEEE Std. 1619.1-2007.\nAmong the new AES-based modes is the NIST approved (see NIST Special Publication SP800-38D defining Galois/Counter Mode (GCM) and Galois Message Authentication Code (GMAC)) GCM mode and IEEE P1619 legacy mode Liskov, Rivest, and Wagner (LRW), that both use Galois Field multiplication for processing 128-bit blocks of data. Besides memory and storage applications, GCM-AES is becoming more widely used in various Internet security protocols and was suggested/submitted as an Internet-draft to the Internet Engineering Task Force (IEFT) to use in the Secure RTP (SRTP) protocol (see Internet-Draft for GCM in Secure RTP (SRTP)), MACsec (see IEEE 802.1AE), Internet Key exchange version 2 (IKEv2), and in the IPsec (see RFC 4106 and RFC 4543).\nA feature of the GCM mode is that the message authentication is performed in parallel with encryption/decryption of the main data payload by applying multiplication in a Galois Field (GF). Multiplications in finite fields have been used for fast (i.e., insecure) message hash computations. To make such computed massage hash values secure, application of the GCM GHASH process adds a pseudorandom vector, a so called \u201cwhitening\u201d vector, at the end. The pseudorandom vector is generated by encrypting a preset value (i.e., Initialization Vector IV) with a secret AES key (i.e., vector W). Use of the GF multiplier for Message Authentication Code (MAC) computation permits higher throughput than the authentication process for computing a conventional MAC. The conventional MAC processes use slower chaining modes, like AES-CBC, or use a separate stand-alone secure hash process from the Secure Hash Algorithm (SHA) family."}
-{"text": "Haematite, having the chemical formula Fe2O3, is one of the most abundant minerals in nature. It exists as iron ore, in other minerals such as bauxite, and is also a component in clay minerals. It is the major component in laeritic soils (red soils found in the tropics). Similarly, manganese oxide, having a formula Mn2O3 is also a very common component in several laeritic soils and also exists as a mineral of manganese in the tropics.\nU.S. Pat. Nos. 5,645,518 and 5,830,815 issued to Wagh et al. on Jul. 8, 1997 and Nov. 3, 1998, respectively, disclose processes for utilizing phosphate ceramics to encapsulate waste. U.S. Pat. No. 5,846,894 issued to Singh et al. on Dec. 8, 1998 discloses a method to produce phosphate bonded structural products from high volume benign wastes. None of these patents provides a method for utilizing the waste materials of iron and manganese.\nU.S. Pat. No. 6,153,809 issued to Singh et al. Nov. 28, 2000 and U.S. patent application Ser. No. 09/751,655 filed Dec. 29, 2000, publication no. U.S. 2002/0123422 to Wagh et al. represent additional development of the use of chemically bonded phosphate ceramics to useful materials. Each of the aforementioned patents, that is U.S. Pat. No. 5,645,518 issued to Wagh et al., U.S. Pat. No. 5,846,894 issued to Singh et al., U.S. Pat. No. 5,830,815 issued to Wagh et al., U.S. Pat. No. 6,153,809 issued to Singh et al., U.S. Pat. No., 6,133,498 issued to Singh et al. and the above-identified publication no. US 2002/0123422 (patent application Ser. No. 09/751,655) is incorporated herein in their entireties.\nThe phosphate ceramics disclosed in the various patents and publication hereinbefore mentioned illustrate a continuing effort to use the chemically bonded phosphate ceramics disclosed therein for a variety of purposes including the encapsulation of hazardous or radioactive waste as seen in the aforementioned publication, as well as the production of low cost structural materials. Accordingly, therefore, a need exists in the art for a low cost structural material which combines with synthetic organic resin based structures, for particular usage in the construction industry. Typically, in warm weather climates, low cost housing may be constructed using styrofoam as a base material onto which is sprayed a concrete-like material as a finish coating to seal the styrofoam base material against the elements and to provide a satisfactory looking structure. Heretofore, the phosphate ceramics disclosed in the above-captioned patents and publication were used as a finish coating in warm temperature climates but have not been satisfactory because the bond between styrofoam and the phosphate ceramics herein above disclosed is physical and peelable such that durable coatings have not been able to be provided with the extant material."}
-{"text": "The invention relates to a tool-driving device that is particularly designed for use in machine tools or in machining units of machining centers, and has at least one machine spindle that is seated to move.\nMachine tools are used especially for material-removal processes, such as boring, milling, turning on a lathe, etc.\nThe tool is inserted into a corresponding tool receptacle that is secured in the work spindle of the relevant machine tool. Various tool receptacles are available.\nDuring the machining process, the work spindles are driven by associated drive apparatuses. Control devices, which can include expanded electronic circuits or execution programs, are provided for controlling the spindle movement, notably its rotation and/or adjustment.\nThe control device establishes the rpm of the spindle within an rpm range. This range is inherently limited. It may be that, particularly for very small tool diameters or for other reasons, rpms outside of the rpm range of the spindle are required.\nIt is the object of the invention to provide a tool-driving device that expands the application range of a machine tool or machining unit, preferably with as little intrusion as possible into the existing machine control.\nThis object is accomplished with a tool-driving device having the features of claim 1.\nThe tool-driving device of the invention has a spindle insert, which can preferably be clamped, fixed against relative rotation, in a machine spindle and can support a tool for machining workpieces. A coupling device serves to secure the spindle insert in the machine spindle. A drive that is supplied by a drive source located outside of the spindle insert, and can be controlled by a control device, is provided for driving the tool. The drive is effected by way of a coupling element that can be connected to the supply lines of the drive. The drive is controlled as a function of the movement of the machine spindle; the tool-driving device is provided with a detection device for detecting this movement.\nFrom the spindle movement, the detection device obtains a signal that characterizes, for example, the rpm, and is used as an input signal for the control device for controlling the drive, and therefore the movement, of the tool. The detection of the rpm requires no access to the machine control, especially if no control signals originating from the machine control are necessary. The control device is separate from the other machine control, and is therefore independent and self-sufficient.\nIf desired, the power supply can be effected by the tapping of the machine control or the drive source of the machine tool. A dedicated drive source can, however, also be provided for the power supply.\nThe tool-driving device permits the increase of the rpm of the machine spindles above and beyond the capabilities of the machine spindle. Unlike in a passive accessory gear, in this instance the additional supply of power in the drive of the tool permits the conversion of an output that exceeds the output of the machine spindle. The maximum torque can be completely retained while the rpm is increased.\nThe spindle insert has a coupling device, e.g., a 7/24 taper shank, which permits a secure, detachable connectionxe2x80x94fixed against relative rotationxe2x80x94with the machine spindle. It also has an essentially cylindrical, one- or multiple-part housing, inside which the drive is disposed.\nIf material-removal operations are to be executed with a rotating tool, the drive is embodied as a rotary drive. A motor, e.g., an electric motor, serves to drive the tool. DC motors, synchronous motors or asynchronous motors can be used for a single- or polyphase alternating current. Hydraulic or pneumatic drives, with which rotational or axial movements of the tool can be attained, can also be used. The motors can be connected to the tool directly, or via a gear in a driving arrangement.\nIn a preferred embodiment of the invention, a receiving apparatus is provided for receiving the tool; the apparatus has a tool spindle, into which the tool is clamped, fixed against relative rotation. The tool spindle preferably has a conical inside shape. The tool spindle is then formed by a rotatably-seated shaft, and projects out of the housing. The shaft is connected to a rotating part of the motor (internal or external rotor) so as to be fixed against relative rotation. The shaft and the tool spindle are preferably embodied to rotate symmetrically relative to an axis of rotation established by the machine spindle. The tool spindle can, however, also support a quick-clamping element, a jaw chuck or the like.\nAt least one slip ring, which is mounted to the outside of the housing and is electrically insulated from it, and can be brought into engagement with an associated sliding contact of the coupling element, is provided for supplying power to the electric motor. When the machine spindle rotates, the sliding contacts slide along the slip rings, thereby assuring the power supply to the drive. Rollers can also be used instead of sliding contacts. The supply can also be effected contactless, e.g., with transformers.\nThe slip rings are preferably disposed on a conical part of the housing whose diameter increases starting from the machine spindle. The slip rings therefore have different diameters. The smallest diameter is larger than that of an arbitrary part of the coupling device. Thus, the spindle insert can be inserted into the machine spindle without altering the position of the contact set. The contact set can then be rigidly secured to the machine tool, in which case it is disposed at a slight incline, corresponding to the incline of the conical housing part. The insertion of the spindle insert produces the contact between the slip rings and the sliding contacts. The contact set can also be seated to be adjusted, and/or can be separate.\nThe safety of the tool-driving device is increased when voltage is only applied to the sliding contacts during the machining process. If the detection device detects rpms that are at least as high as a defined threshold value, preferably 30 rpm, the current supply to the sliding contacts is enabled, for example, by the automatic closure of a switch. The circuit is opened at rpms below the threshold value.\nContactless, magnetic or optical methods are preferred for rpm detection. For example, a metal part connected to rotate with the spindle insert or the machine spindle can serve to induce a short voltage pulse in a stationary coil with each rotation.\nIn an advantageous embodiment, the detection device has a signal generator, particularly a light source, and a signal receiver, particularly a light sensor. The detection device is preferably adjustably mounted to the machine tool, for example to the spindle head that guides the machine spindle. A marking, such as a narrow metal plate, that reflects the light emitted by the light source is secured to the tool coupling or the machine spindle. A signal that is thereby generated, and characterizes the rpm of the machine spindle, e.g., a pulse signal that is proportional thereto, is then transmitted to the control device.\nA circular clamping body having different visual properties from the location where it is to be secured can serve as a marking. The clamping body can have a gap or a recess.\nThe passage of the gap or recess in front of the sensor generates the signal.\nMarkings that effect the generation of a plurality of signals with each rotation can also be provided. In the simplest case, the markings can be equidistantly spaced and provided on, for example, an adhesive strip.\nThe control device utilizes the signals arriving from the detection device to generate a corresponding drive signal for the drive. Hence, the rpm range of the tool can be expanded with the device of the invention. Existing machine tools can therefore be rendered more versatile without its mechanical or electronic components being disturbed.\nThe control device can be integrated into the spindle insert, or accommodated separately. It can also be controlled by programs running on a computer. A console can be provided for the user.\nAt least one supply line for a cooling fluid or compressed air is preferably provided in the tool-driving device for cooling the tool, as is an outward-oriented nozzle, which is preferably pivotable and comprises plastic, for example. At the same time, the nozzle can conduct heat out of the tool-driving device.\nFurther advantageous details about embodiments of the invention ensue from the dependent claims, the drawing and/or the associated description."}
-{"text": "The present invention relates to a hydraulic circuit control system for a construction machine in which an operating system of the construction machine, particularly a control lever device, comprises a joystick device of the type generating an electrical operational signal (electric signal) depending on an input amount upon shift of a control lever, and a flow control valve is controlled with the operational signal for controlling the operation of an actuator.\nIn recent construction machines, particularly in those machines that are employed for various kinds of works because of convenience in use as represented by hydraulic excavators, operability has become increasingly valued in making the machines adaptable for a variety of usages. Stated otherwise, taking a hydraulic excavator as an example, the machine must be able to operate a working device as intended by an operator over a wide range from work in which primary importance is put on the amount of work carried out by the machine, e.g., excavation, to work in which fine adjustment is required in operation, e.g., leveling. To that end, it has been proposed to employ a hydraulic circuit control system in which a control lever device comprises an electric joystick for generating an electrical operational signal depending on an input amount upon shift a control lever, and the operational signal is electrically processed to control a flow control valve with a processed signal. Several known examples of such a control system are as follows.\n(1) Japanese Patent No. 2509311 entitled xe2x80x9cWorking Device Control Method for Construction Machinexe2x80x9d\nThis publication discloses a working device control method for a construction machine comprising a hydraulic control valve (operational valve), which is operated through a controller upon manipulation of an electrical lever, and a pump varying device. Modulation control is performed to absorb shocks caused upon operation of the operational valve and the pump varying device by setting a modulation pattern for rise/fall of a circuit pressure and increase/decrease of a pump delivery rate upon operation of the operational valve to restrict a maximum operating speed of the operational valve (maximum change rate of an operational signal) so that a rate of the rise/fall of the circuit pressure and increase/decrease of the pump delivery rate is gradually changed in multiple stages with a working time, and by operating the operational valve and the pump varying device so as not to move faster than the speeds set by the modulation pattern when the circuit pressure rises and falls at a constant rate with a working time. Furthermore, a cavitation is prevented from occurring upon operation of the pump varying device. This publication also discloses that a plurality of modulation patterns for the operational valve are prepared and one of the patterns is set depending on the working condition automatically or manually with selection by an operator.\n(2) JP,B 7-107279 entitled xe2x80x9cWorking Device Control Method for Construction Machinexe2x80x9d\nThis publication discloses an improvement of the modulation control in the above-mentioned (1). At the time when an electrical lever is manipulated from a shift position on the side in one direction toward the side in an opposite direction in a continuous manner and an operational signal from the electrical lever enters the opposite direction side beyond a dead zone corresponding to a neutral position, the modulation pattern having been effective so far is released and another modulation pattern for the opposite direction side is made effective. The operation of a working device and an operating feeling in the lever-reversed operation are thereby matched with each other.\n(3) JP,A 10-37247 entitled xe2x80x9cOperation Control Device and Operation Control methodxe2x80x9d\nThis publication discloses a hydraulic circuit controller for controlling the operation of a working device of a construction machine through a flow control valve, wherein a maximum change rate of an operational signal for the flow control valve is restrained to be not larger than a setting value, and the operation of the working device is controlled by changing the setting value depending on an input amount upon shift of a control lever.\nMeanwhile, there is also known a hydraulic circuit control system in which an actuator speed is controlled by controlling a delivery rate of a hydraulic pump with an operational signal instead of controlling a flow control valve with the operational signal, and a maximum operating speed of a pump displacement varying mechanism is restrained. Several examples of such a hydraulic circuit control system are as follows.\n(4) JP,B 62-13542 entitled xe2x80x9cController for Hydraulic Circuitxe2x80x9d\nThis publication discloses a hydraulic circuit controller for a closed circuit system wherein an actuator speed is controlled to a speed instructed by an operating device by controlling a delivery rate of a hydraulic pump (position of a pump displacement varying mechanism). When an operating speed of the pump displacement varying mechanism is restrained to be not larger a setting maximum speed, the setting maximum speed is changed depending on an input amount upon shift of a control lever, thereby controlling acceleration/deceleration of an actuator.\n(5) JP,B 62-39295 entitled xe2x80x9cControl System for Hydraulic Circuit Apparatusxe2x80x9d\nThis publication discloses that the controller of the above-mentioned (4) is modified so as to detect a condition of the operating device (control lever) instructing the operation to be stopped or made in the reversed direction, and to set the setting maximum speed larger than that in acceleration.\nThe above-described prior art however has the following problems.\nFirst problem: The setting value for restricting the maximum operating speed of the operational valve (flow control valve) (i.e., the maximum change rate of the operational signal) is not set corresponding to individual operating status, i.e., acceleration, deceleration/stop, and lever-reversed condition. Therefore, the operational valve cannot be always controlled at an optimum maximum change rate adapted for the operating status of a construction machine.\nSecond problem: In the lever-reversed operation, the dead zone in the vicinity of a neutral position of the flow control valve is not appropriately handled or not handled at all. When quickly reversing the control lever, therefore, the actuator undergoes a shock or stalls in the vicinity of the neutral position, causing the operator to feel a pause in the operation.\nThird problem: Since the maximum change speed of the operational valve is just restrained to the fixed modulation pattern regardless of the input amount upon shift of the control lever, an appropriate acceleration/deceleration feeling corresponding to the lever shift amount cannot be provided.\nMore specifically, in Japanese Patent No. 2509311 and JP,B 7-107279, the modulation patterns are set for the maximum operating speed of the operational valve in acceleration and deceleration/stop, and in the lever-reversed operation, the maximum operating speed of the operational valve is restricted in accordance with the modulation pattern for deceleration/stop. However, the lever reversing is performed when it is required to quickly change the moving direction of the working device in the case of, e.g., dropping mud from a bucket, bumping a boom against a vertical surface, or avoiding a risk, and a rapid response is demanded until the working device changes the moving direction. Accordingly, restricting the maximum operating speed of the operational valve in the lever-reversed operation in accordance with the modulation pattern for deceleration/stop cannot be the as providing an optimum maximum operating speed for the lever-reversed operation, and hence cannot change the moving direction of the working device with a good response (first problem).\nAlso, according to JP,B 7-107279, as soon as the operational signal indicates a reversed direction, the modulation control performed so far is ceased and another modulation control adapted for the reversed direction is started for the purpose of improving response in the lever-reversed operation disclosed in Japanese Patent No. 2509311. Taking into account a delay in the operation of the actuator responsive to the operational signal, therefore, the actuator is brought into an uncontrolled state at the moment when the operating direction is changed, which leads to a possibility that a substantial shock may occur until the moving direction of the actuator is completely changed (second problem).\nFurther, in Japanese Patent No. 2509311 and JP,B 7-107279, because the modulation pattern is fixed and the maximum operating speed of the operational valve is always restricted to the fixed modulation pattern regardless of the input amount upon shift of the control lever, an appropriate acceleration/deceleration feeling corresponding to the lever shift amount cannot be provided (third problem). In the case of returning the control lever, for example, when the control lever is manipulated so as to operate the operational valve at a speed higher than that set by the modulation pattern, the maximum operating speed of the operational valve is determined by the fixed modulation pattern regardless of a manner in which the control lever is returned, and therefore cannot be adjusted.\nIn JP,A 10-37247, since the maximum operating speed of the operational valve is not set depending on the operating status of the construction machine, the operational valve cannot be controlled at an optimum maximum change rate adapted for the operating status (first problem), and an appropriate acceleration/deceleration feeling corresponding to the lever shift amount cannot be provided (third problem). Furthermore, no consideration is paid on how to handle the lever-reversed operation (second problem).\nIn JP,B 62-13542 and JP,B 62-39295, the position of the pump displacement varying mechanism is controlled in response to an instruction from the operating device to control the pump delivery rate, thereby controlling the actuator speed. That is to say, these are not intended to control the operation of the working device of the construction machine through the flow control valve. Also, in the system of JP,B 62-39295, a plurality of maximum change rates of the operational signal are set as a function of the operational signal. However, because a control target of the control lever is the pump displacement varying mechanism, no consideration is paid to the dead zone in the vicinity of the neutral position of the flow control valve. Accordingly, if the disclosed arrangement is applied to a hydraulic circuit control system for controlling an actuator speed through a flow control valve, the maximum change rate of an operational signal is restrained in a similar manner even when the flow control valve is within the dead zone in the vicinity of its neutral position, whereby an actuator stalls for a certain period of time, causing the operator to feel a pause in the operation (second problem).\nA first object of the present invention is to provide a hydraulic circuit control system for a construction machine of the type controlling a flow control valve with an electrical operational signal to control the operation of an actuator, the control system being able to control the flow control valve at an optimum maximum change rate in any operating status of acceleration, deceleration/stop, and lever-reversed condition with resulting characteristics cited below:\n(a) in acceleration/deceleration, the machine undergoes a less shock and an operator feels no delay in the operation even with the operator manipulating a control lever quickly;\n(b) in moderate acceleration/deceleration, the actuator is moved as intended by the operator;\n(c) in stop operation, the machine undergoes a less shock and the operator feels no delay in motion toward stop even with the operator manipulating the control lever quickly; and\n(d) in quick lever reversing, the actuator can be rapidly reversed in motion.\nA second object of the present invention is to provide a hydraulic circuit control system for a construction machine, which carries out, in addition to the above, proper processing for a dead zone in the vicinity of a neutral position of the flow control valve in the lever-reversed operation, whereby the machine undergoes a less shock and the operator feels neither a delay in the operation nor a pause in the operation in the vicinity of the neutral position when the control lever is quickly reversed.\nA third object of the present invention is to provide a hydraulic circuit control system for a construction machine, which can give the operator an appropriate feeling in acceleration and deceleration corresponding to an input amount upon shift of the control lever.\n(1) To achieve the above first object, the present invention provides a hydraulic circuit control system for a construction machine comprising a hydraulic actuator for driving a working device, a hydraulic pump driven by a prime mover and producing a pressurized hydraulic fluid, a flow control valve disposed between the hydraulic actuator and the hydraulic pump and controlling a flow rate of the hydraulic fluid, and operational signal generating means for generating an electrical operational signal to instruct a flow rate of the hydraulic fluid flowing through the flow control valve, the system computing a control signal while restraining a change rate of the operational signal to be kept not more than a preset maximum change rate, and controlling the flow control valve in accordance with the computed control signal, wherein the system comprises first determining means for determining the operating status of the construction machine based on the operational signal; and first processing means for setting therein an optimum maximum change rate of the control signal for the flow control valve beforehand for each operating status of the construction machine, determining an optimum maximum change rate adapted for the operating status of the construction machine at that time based on a determination result of the first determining means, and setting the determined optimum maximum change rate as a maximum change rate of the control signal for the flow control valve.\nThus, since the first determining means determines the operating status of the construction machine and first processing means determines an optimum maximum change rate adapted for the operating status of the construction machine at that time based on a determination result of the first determining means and then sets the determined optimum maximum change rate as a maximum change rate of the control signal for the flow control valve, the change rate of the control signal for controlling the flow rate through the flow control valve is restrained to be kept not more than the determined optimum maximum change rate. Therefore, the flow control valve can be controlled at the optimum maximum change rate in any operating status of acceleration, deceleration/stop, and lever-reversed condition with such resulting characteristics as (a) in acceleration/deceleration, the machine undergoes a less shock and an operator feels no delay in the operation even with the operator manipulating a control lever quickly; (b) in moderate acceleration/deceleration, the actuator is moved as intended by the operator; (c) in operation for stop, the machine undergoes a less shock and the operator feels no delay in the motion toward stop even with the operator manipulating the control lever quickly; and (d) in quick lever reversing, the actuator can be rapidly reversed in motion, whereby working efficiency and safety are improved.\n(2) To achieve the above second object, according to the present invention, in the hydraulic circuit control system for a construction machine of the above-mentioned (1), the system further comprises second determining means for determining whether a value of the control signal for the flow control valve is within a neutral zone; and second processing means for computing the control signal in accordance with the operational signal when the value of the control signal for the flow control valve is within the neutral zone, instead of executing the processing to restrain the change rate of the control signal in accordance with the maximum change rate.\nWith those features, proper processing for a dead zone in the vicinity of the neutral position of the flow control valve is executed in the lever-reversed operation so that, when the control lever is quickly reversed, the machine undergoes a less shock and the operation can be performed without causing the operator to feel neither a delay in the operation nor a pause in the operation in the vicinity of the neutral position. As a result, operability in the lever-reversed operation is greatly improved.\n(3) In the above-mentioned (1), preferably, the first determining means determines, based on a state of the operational signal, in which one of acceleration, deceleration/stop, and lever-reversed condition the operating status of the hydraulic excavator is, and the first processing means determines the optimum maximum change rate adapted for the operating status of the construction machine at that time based on the optimum maximum change rate of the control signal set beforehand for each operating status of acceleration, deceleration/stop, or lever-reversed condition.\nWith those features, as with the above-mentioned (1), the flow control valve can be controlled at the optimum maximum change rate in any operating status of acceleration, deceleration/stop, and lever-reversed condition.\n(4) Also, in the above-mentioned (1) or (3), preferably, the first determining means determines the operating status of the construction machine based on the operational signal and a previously outputted control signal for the flow control valve.\nWith that feature, the first determining means can determine the operating status of the construction machine including acceleration, deceleration/stop, and lever-reversed condition.\n(5) To achieve the above third object, according to the present invention, in any one of the above-mentioned (1), (3) and (4), the optimum maximum change rate of the control signal for the flow control valve is set beforehand as a function of the operational signal for each operating status of the construction machine, and the first processing means computes the optimum maximum change rate based on the function of the operational signal corresponding to the operating status determined by the first determining means and the operational signal at that time.\nWith those features, the optimum maximum change rate of the control signal is set depending the value of the operational signal, and hence an appropriate feeling in acceleration and deceleration corresponding to the input amount upon shift of the control lever can be provided.\n(6) In any one of the above-mentioned (1), (3) and (4), preferably, the optimum maximum change rate of the control signal for the flow control valve is set beforehand as a function of the operational signal or a function of the previously outputted control signal for the flow control valve for each operating status of the construction machine, and the first processing means computes the optimum maximum change rate based on the function of the operational signal corresponding to the operating status determined by the first determining means or the function of the previously outputted control signal for the flow control valve and the operational signal at that time or the previously outputted control signal for the flow control valve.\nWith those features, the optimum maximum change rate of the control signal is set depending both the value of the operational signal and the previously outputted control signal, and hence an appropriate feeling in acceleration and deceleration corresponding to the input amount upon shift of the control lever can be provided."}
-{"text": "The present invention relates to an apparatus and method for the solidification of sludges by kneading sludges in soft ground layers with a hardener and solidifying the sludges.\nAs the conventional method for solidifying sludges deposited and accumulated in bottom portions of harbors, bays, rivers and lakes, there can be mentioned a method in which a hardener is added to sludge in situ and the hardener-incorporated sludge is kneaded. The reason why the sludge is solidified in situ in the deposited and accumulated state is that the amount of the water contained in the sludge is held to a minimum and the solidification treatment can be performed conveniently. If the sludge is dug out and placed on the land for solidification, the water content in the sludge is increased greatly compared with the sludge deposited naturally.\nThis known method, however, is defective in that when the deposited sludge is kneaded and agitated by a kneading machine or the sludge is dredged, sea water or the like is contaminated in a broad region to cause secondary pollutions such as the generation of bad odors. Moreover, when the sludge is agitated, water or untreated sludge flows from neighboring sludge layers into the sludge being treated, and therefore, a large quantity of a hardener must be added. In this case, the hardener supplied in such a large amount readily flows into neighboring sludge layers and is wasted. Agitation of the sludge or the like is performed by agitating blades of the kneading machine. Accordingly, the sludge being treated is not completely separated from the sludge present on or around the outer periphery of the rotation locus of the agitating blades and therefore, uniform kneading of the sludge and hardener is not attained, and it is difficult to perform the solidification treatment in good kneading conditions."}
-{"text": "1. Field of the Invention\nThe present invention relates to a technology for realizing low damage sputtering regardless of materials (such as, an inorganic or organic material) in a surface analysis method, and relates to a technology for realizing improvement of sensitivity by improving secondary ion yield in a secondary ion mass spectroscopy method.\n2. Description of the Related Art\nAn ion source for surface analysis that can perform sputtering without any damage to a target sample has not yet been developed. In the surface analysis, argon ion (Ar+) is the most common ion species for sputtering, but it is known that the occurrence of damage due to the sputtering is high.\nIn addition, in a secondary ion mass spectroscopy method (SIMS) as one of surface analysis methods, a primary ion beam that has been used so far is a noble gas ion or a metal ion (Cs+, Ar+, Ga+, Au+, or the like). Some of them can be reduced to a small beam in the order of several tens nanometer, but the occurrence of large damage to a sample is a common drawback.\nIn addition, if these ions are used as a primary ion source, secondary ion yield is very low, and secondary ion generation efficiency is low. Therefore, in order to overcome the drawback of the SIMS using them as the primary beam, a cluster ion SIMS has been developed. A beam source thereof is Au3++, Bi3++, or the like. By using the cluster ion (Au3++, Bi3++, or the like) consisting several atoms, desorption efficiency of the secondary ions is significantly increased in a non-linear manner. Such result is due to the generation of ablation.\nOn the other hand, because a target sample surface and its vicinity are significantly damaged, application of the conventional system to a biological material is difficult; and nondestructive observation of molecule ions is difficult; specifically, the sample receives large fragmentation, and a surface of the sample is decomposed and polymerized.\nA cluster ion source of C60+ ion is commercialized; and hence, a low damage sputtering technology is realized though in a limited manner. Further, the desorption efficiency is further increased in the SIMS using the C60+ ion source as the primary ion source. However, the following phenomena are caused: (1) an inorganic material is contaminated with a carbon component derived from C60; (2) craters are generated in a surface of the material so that surface destruction occurs; (3) a biological sample or the like is significantly damaged; and (4) the secondary ion yield is low in the SIMS, and when the beam diameter is decreased, ionic strength is weakened so that utility value as the SIMS is deteriorated (particularly in an organic material). Refer to Japanese Patent Application Laid-open No. 2005-134170, Journal of Physical Chemistry B, 108, pp 7831-7838, and Applied Surface Science 231-232, pp 936-939, FIG. 4.\nThere is a surface analysis method utilizing a gas cluster ion beam (GCIB) that has been recently popular, in which noble gas (such as argon (Ar)) is ejected in vacuum to form a jet stream, gas temperature is decreased, and neutral clusters having an n value of Arn+ of a few thousands to a few tens of thousands are formed and ionized to generate Arn+, which is accelerated to impact the sample.\nWith this method, depth profile analysis with low-damage sputtering for an organic material (such as a polymer) is confirmed to be effective and is commercialized. However, for an inorganic material (such as a ceramic material) that is relatively hard, the sputtering speed is extremely slow so that it is not practical. Therefore, a range of the sample types to be analyzed is inevitably limited to mainly organic industrial materials.\nIn addition, when the GCIB is used as the primary ion source in the secondary ion mass spectroscopy method, it is known that the secondary ion yield thereof is low; and hence, it is not practical when used for improving sensitivity in the secondary ion mass spectroscopy method. Refer to Japanese Patent Application Laid-open No. Hei 04-354865, Japanese Patent Application Laid-open No. 2008-116363, and Analytical Chemistry, 2011, 83(10), pp 3793-3800, FIG. 7.\nIn addition, an ion beam technology using a charged droplet method has been developed. In this method, a capillary is disposed in the atmosphere, solvent is supplied through inside of the capillary, and an extraction electrode that is applied with a high voltage negative with respect to the capillary is disposed in front of the capillary so as to generate ions in the atmosphere.\nA vacuum chamber is separated into several steps from low vacuum side to high vacuum side with small diameter orifices. The ions are made to pass through the orifices and are transported to vacuum atmosphere so as to be used as ion beam. In this case, the cluster ions generated in the atmosphere inevitably collide with gas molecules in the atmosphere so that many ions are scattered. Therefore, the amount of ions that are actually transported to the vacuum side and can be effectively used is small; and in addition, downsizing of the cluster ion (fission of the cluster) also occurs due to vaporization in the atmosphere side.\nIn addition, to use the ion beam, it is necessary to apply a high voltage, which is positive with respect to the ground potential, to the capillary as a source, and it is also necessary to apply a high voltage to parts for lens effect or the like in a low vacuum region during the ion transportation process. Therefore, discharge phenomenon tends to occur in various parts. Consequently, it becomes difficult to stably obtain the ion beam, and it is also difficult to decrease the beam size to be small.\nOn the other hand, a differential pumping system for evacuating the separated vacuum chamber also becomes large in scale which causes difficulty when in use. Refer to Japanese Patent Application Laid-open No. 2011-141199.\nConsequently, a practical ion source that can support various types in etching layer-by-layer without damaging a surface of the sample after irradiation has not been developed yet, and an ion source succeeding in dramatic improvement of sensitivity in the secondary ion mass spectroscopy method has also not yet been developed.\nA charged droplet ion source of the related art is described below. In FIG. 5, a charged droplet ion source 701 includes a vacuum chamber 710.\nThe vacuum chamber 710 is connected to first and second vacuum evacuating devices 729a and 729b so that the inside of the vacuum chamber 710 can be evacuated.\nAn extracting electrode 721 is provided with a small hole (orifice) so that gas flows in the vacuum chamber 710 through the extracting electrode 721 when the inside of the vacuum chamber 710 is evacuated. First, the inside of the vacuum chamber 710 is evacuated by the first and second vacuum evacuating devices 729a and 729b. \nAn emission tube (capillary) 703 is disposed outside the vacuum chamber 710.\nThe distal end of the emission tube 703 is directed towards the small hole of the extracting electrode 721; and a base part thereof on the opposite side is connected to a liquid supply pipe 743. The liquid supply pipe 743 is connected to an ionization liquid supply device 705.\nThe ionization liquid supply device 705 includes a liquid storing portion 732 and a liquid feeding pump 731. The ionization liquid stored in the liquid storing portion 732 is supplied to the base part of the emission tube 703 through the liquid supply pipe 743 by the liquid feeding pump 731, passes a thin tube in the emission tube 703, and is emitted to the outside of the emission tube 703 from an emission opening 735 at the distal end of the emission tube 703. The emission tube 703 is surrounded by an outer cylinder 707. When carrier gas (here, nitrogen gas) is supplied from a carrier gas source 708 to the inside of the outer cylinder 707, the gas is released from a distal end opening 736 of the outer cylinder 707.\nThe emission opening 735 is disposed between the distal end opening 736 of the outer cylinder and the small hole of the extracting electrode 721. Around the emission opening 735, there is formed a flow of the carrier gas from an upstream side as the base side of the emission tube 703 to a downstream side on which the extracting electrode 721 is located with the small hole.\nAn extraction power supply 728 is disposed outside the vacuum chamber 710.\nIn a state where the carrier gas supplied from the carrier gas source 708 is released from the distal end opening 736, the liquid feeding pump 731 supplies the ionization liquid to the emission opening 735, the extraction power supply 728 applies a voltage between the emission tube 703 (made of a metal here) and the extracting electrode 721 so that an electric field thereof extracts droplet cluster ions charged with a positive charge from the ionization liquid positioned in the emission opening 735. Then, the cluster ions pass through the small hole of the extracting electrode 721 and enter the inside of the vacuum chamber 710.\nOn the downstream side of the extracting electrode 721, there are disposed accelerating electrodes 722 and 723 with small holes and transport lens electrodes 724 and 725. When voltages are applied to the electrodes 722 to 725, the droplet cluster ions entering the inside of the vacuum chamber 710 pass through holes formed in the electrodes 722 to 725 so as to be a droplet cluster ion beam, and further propagates toward the downstream side.\nA size of an initial droplet cluster ion generated in the atmosphere is approximately 100 nm in diameter. However, the droplet cluster ion generated in the atmosphere is downsized due to Rayleigh fission that occurs when Coulomb repulsion of itself exceeds surface tension of the droplet. Further, the droplet cluster ions inevitably collide with gas molecules in the atmosphere so that many ions are scattered. Therefore, only a small amount of the droplet cluster ions can enter the inside of the vacuum chamber 710, and the size of the droplet cluster ion is decreased to be smaller than that of initially generated one.\nIn addition, for use as the droplet cluster ion beam, it is necessary to apply a positive high voltage with respect to the ground potential to the emission tube 703 as the generation source. Further, it is also necessary to apply high voltages to the extracting electrode 721, the first accelerating electrode 722, and the transport lens electrode 724 disposed in the low vacuum environments in the vacuum chamber 710. Therefore, an arcing phenomenon is apt to occur in the vacuum chamber 710, and hence it is difficult to obtain the droplet cluster ion beam.\nIn addition, it is necessary to separate the atmosphere outside the vacuum chamber 710 from the inside space of the vacuum chamber 710, both of which are connected to each other through the small hole of the extracting electrode 721. Therefore, the first and second vacuum evacuating devices 729a and 729b for evacuating the inside space of the vacuum chamber 710 are required to be large ones; and hence, difficulty arises when they are used in that they occupy large areas and in terms of cost.\nConsequently, in the ion source on the conventional technology, disposing the emission opening of the emission tube in the atmosphere so that the droplet cluster ion beam is generated in the atmosphere provides small amount of the droplet cluster ions that can be actually used. Hence, the conventional technology is of little practical use."}
-{"text": "The present invention relates to intermediate molecular weight shaped polyethylene articles such as polyethylene fibers exhibiting relatively high tenacity, modulus and toughness, and to products made therefrom. The polyethylene article is made by a process which includes the step of stretching a solution of polyethylene dissolved in a solvent at a stretch ratio of at least about 3:1.\nPolyethylene fibers, films and tapes are old in the art. An early patent on this subject appeared in 1937 (G.B. No. 472,051). However, until recently, the tensile properties of such products have been generally unremarkable as compared to competitive materials, such as the polyamides and polyethylene terephthalate. Recently, several methods have been discovered for preparing continuous low and intermediate molecular weight polyethylene fibers of moderate tensile properties. Processes for the production of relatively low molecular weight fibers (a maximum weight average molecular weight, Mw, of about 200,000 or less) have been described in U.S. Pat. Nos. 4,276,348 and 4,228,118 to Wu and Black, U.S. Pat. Nos. 3,962,205, 4,254,072, 4,287,149 and 4,415,522 to Ward and Cappaccio, and U.S. Pat. No. 3,048,465 to Jurgeleit. U.S. Pat. No. 4,268,470 to Cappaccio and Ward describes a process for producing intermediate molecular weight polyolefin fibers (minimum molecular weight of about 300,000).\nThe preparation of high strength, high modulus polyolefin fibers by solution spinning has been described in numerous recent publications and patents. German Off. No. 3,004,699 to Smith et al. (Aug. 21, 1980) describes a process in which polyethylene is first dissolved in a volatile solvent, the solution is spun and cooled to form a gel filament, and, finally, the gel filament is simultaneously stretched and dried to form the desired fiber. U.K. Patent Application No. 2,051,667 to P. Smith and P. J. Lemstra (Jan. 21, 1981) discloses a process in which a solution of a polymer is spun and the filaments are drawn at a stretch ratio which is related to the polymer molecular weight, at a drawing temperature such that at the draw ratio used, the modulus of the filaments is at least 20 GPa (the application notes that to obtain the high modulus values required, drawing must be performed below the melting point of the polyethylene; in general, at most 135.degree. C.). Kalb and Pennings in Polymer Bulletin, Volume 1, pp. 879-80 (1979), J. Mat. Sci., Vol. 15, pp. 2584-90 (1980) and Smook et al. in Polymer Mol., Vol 2, pp. 775-83 (1980) describe a process in which the polyethylene is dissolved in a non-volatile solvent (paraffin oil), the solution is cooled to room temperature to form a gel which is cut into pieces, fed to an extruder and spun into a gel filament, the gel filament being extracted with hexane to remove the parafin oil, vacuum dried and stretched to form the desired fiber.\nMost recently, ultra high molecular weight fibers have been disclosed. U.S. Pat. No. 4,413,110 to Kavesh and Prevorsek describes a solution spun fiber of from 500,000 molecular weight to about 8,000,000 molecular weight which exhibits exceptional modulus and tenacity. U.S. Pat. Nos. 4,430,383 and 4,422,993 to Smith and Lemstra also describe a solution spun and drawn fibers having a minimum molecular weight of about 800,000. U.S. Pat. No. 4,436,689 to Smith, Lemstra, Kirschbaum and Pijers describes solution spun filaments of molecular weight greater than 400,000 (and an Mw/Mn<5). In addition, U.S. Pat. No. 4,268,470 to Ward and Cappacio also discloses high molecular weight polyolefin fibers.\nIn general, the known processes for forming polyethylene and other polyolefin fibers may be observed as belonging in one of two groups: those which describe fibers of low average molecular weight (200,000 or less) and those which describe fibers of high average molecular weight (800,000 or more). Between the two groups, there is a molecular weight range which has not been accessible to the prior art methods for preparing fibers of high tensile properties.\nThere are advantages to the molecular weight ranges thus far mastered. Lower molecular weight polymers are generally synthesized and processed into fibers more easily and economically than high molecular weight fibers. On the other hand, fibers spun from high molecular weight polymers may possess high tensile properties, low creep, and high melting point. A need exists for fibers and methods which bridge this gap, combining good economy with moderate to high tensile properties. Surprisingly, our process makes it possible to accomplish these results."}
-{"text": "1. Field of the Invention\nThe present invention relates to high density memory devices based on phase change memory materials, including chalcogenide based materials and on other programmable resistance materials, and methods for manufacturing such devices.\n2. Description of Related Art\nPhase change based memory materials, like chalcogenide based materials and similar materials, can be caused to change phase between an amorphous state and a crystalline state by application of electrical current at levels suitable for implementation in integrated circuits. The generally amorphous state is characterized by higher electrical resistivity than the generally crystalline state, which can be readily sensed to indicate data. These properties have generated interest in using programmable resistance material to form nonvolatile memory circuits, which can be read and written with random access.\nThe change from the amorphous to the crystalline state is generally a lower current operation. The change from crystalline to amorphous, referred to as reset herein, is generally a higher current operation, which includes a short high current density pulse to melt or breakdown the crystalline structure, after which the phase change material cools quickly, quenching the molten phase change material and allowing at least a portion of the phase change material to stabilize in the amorphous state.\nThe magnitude of the current needed for reset can be reduced by reducing the size of the phase change material element in the cell and/or the contact area between electrodes and the phase change material, so that higher current densities are achieved with small absolute current values through the phase change material.\nOne approach to reducing the size of the phase change element in a memory cell is to form small phase change elements by etching a layer of phase change material. However, reducing the size of the phase change element by etching can result in damage to the phase change material due to non-uniform reactivity with the etchants which can cause the formation of voids, compositional and bonding variations, and the formation of nonvolatile by-products. This damage can result in variations in shape and uniformity of the phase change elements across an array of memory cells, resulting in electrical and mechanical performance issues for the cell.\nAdditionally, it is desirable to reduce the cross-sectional area or footprint of individual memory cells in an array of memory cells in order to achieve higher density memory devices. However, traditional field effect transistor access devices are horizontal structures having a horizontally oriented gate overlying a horizontally oriented channel region, resulting in the field effect transistors having a relatively large cross-sectional area which limits the density of the array. Attempts at reducing the cross-sectional area of horizontally oriented field effect transistors can result in issues in obtaining the current needed to induce phase change because of the relatively low current drive of field effect transistors.\nThus, memory devices including both vertically and horizontally oriented field effect transistors have been proposed. See, for example, U.S. Pat. No. 7,116,593. However, the integration of both vertically and horizontally oriented field effect transistors can be difficult and increase the complexity of designs and manufacturing processes. Thus, issues that devices having both vertically and horizontally oriented field effect transistors need to address include cost and simplicity of manufacturing.\nAlthough bipolar junction transistors and diodes can provide a larger current drive than field effect transistors, it can be difficult to control the current in the memory cell using a bipolar junction transistor or a diode adequately enough to allow for multi-bit operation. Additionally, the integration of bipolar junction transistors with CMOS periphery circuitry is difficult and may result in highly complex designs and manufacturing processes.\nIt is therefore desirable to provide both vertically and horizontally oriented field effect transistors on the same substrate that are readily manufactured for use in high-density memory devices, as well as in other devices that may have a need for both types of transistors on one chip. It is also desirable to provide memory devices providing the current necessary to induce phase change, as well as addressing the etching damage problems described above."}
-{"text": "1. Field of the Invention\nThe present invention relates to virtual reality displays.\n2. Related Art\nPrior methods of virtual reality display systems deployed head or helmet mounted display that placed a viewing screen directly in front of the user's eyes and recorded the movement of the users head to determine what should be shown on the display. Thus, when the head turned to one side, the display was refreshed to show what was in the virtual world in the direction they turned their head."}
-{"text": "1. Technical Field\nThe present invention relates to an improved information-retrieval apparatus. In particular, the present invention relates to an improved digitally-based information-retrieval apparatus. More particularly, the present invention relates to magnetic storage media, such as digital recording tapes. Still more particularly, the present invention relates to magnetic recording heads for writing and reading data to digital recording tapes.\n2. Description of the Related Art\nVarious magnetic recording techniques exist for recording data to and from magnetic storage media, such as magnetic tape. Magnetic tapes are used for data storage in computer systems requiring data removability, low-cost data storage, high data-rate capability, high volumetric efficiency and reusability. The constantly increasing operational speeds of digital computers are creating a demand for corresponding increases in the data storage capacities of magnetic tape recording and reproducing systems, while maintaining the special requirements of high speed digital tape systems.\nTape recording and reproducing systems for use as computer data storage devices are required to provide high data transfer rates and to perform a read check on all written data. To satisfy these requirements, conventional tape systems typically employ methods of recording known as linear recording, in which the tracks of data lie parallel to each other and to the edge of the tape. Linear recording techniques offer high data transfer rates. However, it is desirable to obtain even higher data densities while retaining the advantages of such recording techniques.\nDigital linear tape (DLT) is a magnetic linear tape medium that is increasingly being utilized as a medium for data storage. DLT is a magnetic storage medium used to back up data, typically in computer systems. DLT allows for the rapid transfer of data, in comparison to other tape storage technologies. For example, various forms of magnetic read/write heads can be utilized in association with servo mechanisms to read and write data to and from a track of a particular DLT.\nBecause DLT is currently being utilized as an important tool for data storage, it is desirable to increase the recording density, thus allowing for the faster and more efficient retrieval and writing of data. One method of increasing this storage density involves azimuth recording. The term \"azimuth\" refers to the horizontal angular distance from a particular reference direction. The use of the word \"azimuth\" in \"azimuth recording\" thus suggests a form of angular recording.\nAzimuth recording involves the use of a rotating recording head, such that data tracks on a tape may be recorded at different angles with respect to the edge of the tape. Azimuth recording results in a recorded track pattern in which the magnetization directions of adjacent data tracks lie at different azimuth angle to each other. To date, most recording systems have relied strictly on magnetic heads which contain read/write elements but which record only vertically, thus not allowing for angular or \"azimuthal\" recording of data. One of the principal advantages of azimuth recording over non-azimuth recording is that azimuth recording promotes very high data track packing. Azimuth recording provides much denser track packing than regular track packing spacing because regular track packing spacing typically requires gaps between tracks.\nThose systems which do attempt to implement azimuth recording techniques are faced with the challenge of providing fine positioning servo tracking. Servo tracking techniques have been developed to reduce the effects of tracking error and thus improve the data capacity of tape systems. Known servo techniques vary widely, but most involve methods of dynamically moving a read head gap to continually reposition it over a written servo track. The movement of a servo read head gap compensates for lateral tape motion during a read. However, lateral tape motion during writing is not controlled with respect to the write head gap. Thus, the distance between tracks is still limited to the magnitude of the lateral tape motion in order to avoid over-writing previously written tracks.\nBased on the foregoing it can be appreciated that a need exists for an improved azimuth recording system which does not encounter problems associated fine positioning servo tracking. A need also exists for an inexpensive and easy to implement apparatus and method which provides fine positioning servo tracking. It is believed that the apparatus and method presented herein solves these problems."}
-{"text": "For the purpose of conveying loose materials such as bulk materials, that is, rock/stones, mineral resources, excavation material, agricultural products, et cetera, use has long been made of troughed conveyor belts, which receive the conveyable material at a receiving location on their carrying side and discharge the same at a discharge location. Since the conveyable material is open to the environment as it is being transported, contaminants and environmental weathering influences can act on the conveyable material, and the latter can pollute the environment and also pose a risk to the environment. It is also the case, on account of their configuration, that troughed conveyor belts can be used to realize curves and gradients only to a limited extent. It is thus not usually possible, in conventional belt systems, to exceed an angle of inclination of 20\u00b0 in gradient. If this is the limit of feasibility, it is necessary to connect a plurality of inter alia specific conveying belts with transfer locations. This increases the complexity for, and therefore the costs of, the conveyor system to a considerable extent.\nIn order to eliminate these disadvantages, conveyor belts which are closed during operation and are referred to as tube belts, tubular conveyor belts, pipe belts or mega pipes, were developed in the 1980s. The tube belts are rolled together between the receiving location and discharge location to give a closed tube, by virtue of the outer belt flanks overlapping and thus fully enclosing the conveyable material. This means that the conveyable material in the tube belt and the environment are completely separated from one another, since the tube belt remains closed over the conveying route. It is only for the purposes of receiving and discharging the transportable material that a tube belt widens and assumes the form of a conventional troughed conveyor belt. This rules out contamination of the bulk material along the conveying route and the associated environmental pollution. It is also the case that the conveyable material cannot be influenced by the environment during transportation. Further essential advantages of the tube belts in relation to the conventional troughed conveyor belts reside in the possibility of realizing very narrow three-dimensional curves and in the relatively high angles of inclination of up to 35\u00b0 in gradient, which means that complicated three-dimensional curved routes can be realized by a single system. Since tube belts usually have a smooth surface on their carrying side, the angles of inclination are nevertheless limited to a gradient of up to 35\u00b0, depending on the bulk material properties.\nIn order to eliminate these disadvantages, conveyor belts which are closed during operation and are known as SICON\u00aeconveyor belts or pocket (conveyor) belts have also been developed. A pocket conveyor belt comprises two textile-reinforced profiles each with a steel cable vulcanized therein as a tension member. The profiles run over the sets of rollers and carry the pocket which accommodates the conveyable material. This droplet-shaped pocket consists of highly flexible rubber and is connected to the profiles by means of hot vulcanization. The profiles are arranged one above the other during transportation, and the belt is therefore closed off in a dust-tight manner. The belt is carried, and guided, by specific sets of rollers which, for the closed state of the belt, comprise a carrying roller and a guide roller. Further sets of rollers, each comprising a carrying roller and one to three guide rollers, are available for loading and unloading the belt and for curves and gradients.\nIn a manner similar to tube belts, the essential advantages of the pocket conveyor belt in relation to the conventional troughed conveyor belts reside in the possibilities of realizing very narrow three-dimensional curves and in the relatively high angles of inclination of up to 35\u00b0; in the case of conventional belt systems, the angle of inclination cannot usually exceed 20\u00b0. This makes it possible to realize complicated three-dimensional curved routes by a single system, without any transfer locations on the conveying route. In addition, the material in the pocket conveyor belt and the environment are completely separated from one another, since the pocket conveyor belt remains closed over the conveying route. For loading purposes, the pocket conveyor belt is opened with the aid of a specific set of rollers for opening and closing the belt. The belt can be unloaded at an overhead discharge point or an S-shaped discharge station. At the S-shaped discharge station, it is possible optionally for the belt to be emptied or for the conveyable material to be poured into the belt again.\nA pocket conveyor belt differs from the conventional tube belt not just in construction, but also in functioning and areas of application. It is thus possible for a pocket conveyor belt, depending on the profile size, to negotiate radii of 0.6 m or 1.0 m, which cannot be realized by a conventional tube belt. The minimum curve radius which can be realized by a tube belt is approximately 30 m. In contrast to the tube belt, the conveyor length, the conveyor cross section and the associated conveyor capacity and maximum possible material particle size of a pocket conveyor belt are very limited. All of this predestines a pocket conveyor belt for an \u201cin-plant closed\u201d transport of industrial bulk materials, while a tube belt is considered in practice to be more akin to an \u201cout-plant closed\u201d conveying principle for the entire range of particle sizes.\nFor a number of application cases, the advantages of the tube belts or the pocket conveyor belts and the steep conveyor belts are required at the same time, that is, a tube belt or pocket conveyor belt which can be used even at angles of inclination above 35\u00b0.\nU.S. Pat. No. 6,170,646, GB 1197700, U.S. Pat. No. 5,351,810, JP 480 48 385 U, JP 580 83 314 U, United States patent application publication 2012/0000751 A1, FR 14 968 97 A, GB 88 76 98 A, JP 582 16 803 A, U.S. Pat. No. 3,392,817 A and WO2005/085101 A1 disclose a number of technical solutions in this respect for increasing the angle of inclination of tube belts and pocket belts by differently shaped profiles having been applied to the carrying-side cover panel of a tube belt or pocket belt. The core idea of these approaches has been in each case, for elastic rubber or plastic-material strips connected to the conveyor belt to be fitted transversely to the longitudinal direction of the conveyor belt and to be offset at certain intervals from one another in the longitudinal direction. It is possible here for the transverse strips to span both the entire belt width and just part of the belt width. It can be established from these documents that the transverse strips may be configured both in a continuous state, in the form of ribs or wave-like strips, and in a divided state, for example, at right angles, in sawtooth form or in trapezoidal form. The divided transverse strips here are configured such that, when the tube belt or pocket belt is deformed in tube or pocket form, the flanks of the strip butt more or less against one another or overlap and thus form partition walls spaced apart in the longitudinal direction. Depending on the height, that is, radial formation, of the transverse strips, the conveyable material is retained in a force-fitting and form-fitting manner during transportation, and it is therefore possible to prevent the conveyable material from sliding back in the conveyor belt and thus to realize relatively large gradients. In the case of the transverse strips being virtually closed, it is even possible to realize vertical conveying directions, wherein purely form-fitting force transmission takes place.\nIt is a disadvantage of the above-described tube belts or pocket belts that they involve very high outlay, and are therefore expensive to produce. It is thus necessary for the transverse strips, on account of their size, in particular their radial extent, to be produced in the form of separate elements and to be applied subsequently to the conveyor belt for example by means of adhesive bonding, that is, by cold vulcanization. This requires the further operating steps of the transverse strips being separately produced and subsequently installed on the conveyor belt. Single-piece production of conveyor belts with transverse strips, that is, simultaneously with the vulcanization of the conveyor belt, is ruled out in production terms on account of the size of the transverse strips. It may also be necessary for the transverse strips to be installed on the conveyor belt for the first time at the site of use of the conveyor belt, so that there is no increase in the volume of the conveyor belt for transportation purposes. Furthermore, the adhesive-bonding locations constitute a weak point which, over time, will fail sooner than other constituent parts of the conveyor belt.\nIt is also disadvantageous that, if the known tube belts or pocket belts are suitable for relatively large angles of inclination, that is, above 35\u00b0 in gradient, the conveyable material is retained in a form-fitting manner by the transverse strips and the latter are subjected to corresponding loading. This requires a corresponding stable and radial formation of the transverse strips with higher material and production costs than in the case of flatter profilings, although the latter do not allow such gradients. It is also the case that the higher transverse strips increase the transportation costs of the conveyor belts, because the latter cannot be wound as tightly for transportation purposes, that is, less belt length per rolled together belt drum can be transported in one journey. At the same time, this means that the pieces of belt which can be transported per drum in one journey are shorter, and there is therefore an increase in the outlay for installing the endlessly closed conveyor belts in the conveyor-belt system. As is also the case with conventional tube belts, the pocket conveyor belts have a smooth surface, as a result of which it is possible to realize the angles of inclination of up to 35\u00b0, depending on the bulk-material properties."}
-{"text": "1. Field of the Invention\nThe present invention relates to a method and related device for controlling operation of a portable electronic device, and more particularly, to a method and related device for determining whether a lid of a portable electronic device is open using a gravitational acceleration sensor and correspondingly determining the operation of the portable electronic device.\n2. Description of the Prior Art\nA laptop (i.e., a notebook computer) has several advantages, such as a small-sized volume, lightweight, and convenient for carrying due to its portability. These properties allow a user to work in any location. A small, thin, and light notebook computer provides the user with powerful computation abilities and document or multimedia processing functions anywhere and anytime, and thereby the work location of the user is not limited.\nPlease refer to FIG. 1. FIG. 1 is a schematic diagram of a notebook computer system 10 according to the prior art. Generally speaking, the notebook computer system 10 is composed of a lid 100 and a base 102. A hinge 104 connects the lid 100 and the base 102. The lid 100 comprises a screen, a camera, etc. The base 102 comprises a keyboard, a touchpad, a power switch, a host, an expanding interface, and so on. When using the notebook computer system 10, the user has to turn on the power of the host and adjusts a display angle of the screen of the lid 100 to a specific angle. In order to save power, a switch installed in the notebook computer system 10 can switch ON/OFF states of the screen and operation of the host according to an opening angle of the lid 100. For example, when the user doesn't need to use the notebook computer system 10 after turning on the notebook computer system 10, the user can close the lid 100 to make an angle, between the lid 100 and the base 102, smaller than a specific value, so that the notebook computer system 10 turns off the screen and operates in a sleep mode.\nAdjusting the angle between the lid 100 and the base 102, the user can save power and timely switch the operation of the notebook computer system 10. Therefore, it is considerably important to precisely detect the angle between the lid 100 and the base 102. In the prior art, there are many ways to detect the opening angle of the lid 100 and one of these is using a mechanic switch connected to the hinge 104. That is, turning off the screen and executing related operations, e.g. operating in the sleep mode when a rotating angle of the hinge 104 is smaller than a specific angle. However, the assembly of the mechanic switch is difficult and the mechanic switch may weary or malfunction by time, and finally, the reliability of the mechanic switch is decreased.\nIn addition, a magnetic sensor, such as a Hall sensor or a magnetic reluctance sensor, is used in the prior art. The notebook computer system 10 receives a distance from the lid 100 to the base 102 for determining the angle between the lid 100 and the base 102. For example, the Hall sensor can sense magnetic pole and magnetic force. Therefore, by installing a magnet in the lid 100 and a Hall sensor in the base 102, the notebook computer system 10 can determine the distance between the lid 100 and the base 102 so as to determine the angle between the lid 100 and the base 102. However, it is necessary to take the sensibility of the Hall sensor and magnetic flux of the magnet into account to meet demands when installing the magnet and the sensor. Besides, the magnetic reluctance sensor is difficult to design because of its high sensibility and narrow linear range."}
-{"text": "This invention relates to preparing fluorinated electrets.\nThe filtration properties of nonwoven polymeric fibrous webs can be improved by transforming the web into an electret, i.e., a dielectric material exhibiting a quasi-permanent electrical charge. Electrets are effective in enhancing particle capture in aerosol filters. Electrets are useful in a variety of devices including, e.g., air filters, face masks, and respirators, and as electrostatic elements in electro-acoustic devices such as microphones, headphones, and electrostatic recorders.\nElectrets are currently produced by a variety of methods including direct current (xe2x80x9cDCxe2x80x9d) corona charging (see, e.g., U.S. Pat. Re. 30,782 (van Turnhout)), and hydrocharging (see, e.g., U.S. Pat. No. 5,496,507 (Angadjivand et al.)), and can be improved by incorporating fluorochemicals into the melt used to produce the fibers of some electrets (see, e.g., U.S. Pat. No. 5,025,052 (Crater et al.)).\nMany of the particles and contaminants with which electret filters come into contact interfere with the filtering capabilities of the webs. Liquid aerosols, for example, particularly oily aerosols, tend to cause electret filters to lose their electret enhanced filtering efficiency (see, e.g., U.S. Pat. No. 5,411,576 (Jones et al.)).\nNumerous methods have been developed to compensate for loss of filtering efficiency. One method includes increasing the amount of the nonwoven polymeric web in the electret filter by adding layers of web or increasing the thickness of the electret filter. The additional web, however, increases the breathing resistance of the electret filter, adds weight and bulk to the electret filter, and increases the cost of the electret filter. Another method for improving an electret filter\"\"s resistance to oily aerosols includes forming the electret filter from resins that include melt processable fluorochemical additives such as fluorochemical oxazolidinones, fluorochemical piperazines, and perfluorinated alkanes. (See, e.g., U.S. Pat. No. 5,025,052 (Crater et al.)). The fluorochemicals should be melt processable, i.e., suffer substantially no degradation under the melt processing conditions used to form the microfibers that are used in the fibrous webs of some electrets. (See, e.g., WO 97/07272 (Minnesota Mining and Manufacturing)).\nIn one aspect, the invention features an electret that includes a surface modified polymeric article having surface fluorination produced by fluorinating a polymeric article. In one embodiment, the article includes at least about 45 atomic % fluorine as detected by ESCA. In another embodiment, the article includes a CF3:CF2 ratio of at least about 0.25 as determined according to the Method for Determining CF3:CF2. In other embodiments, the article includes a CF3:CF2 ratio of at least about 0.45 as determined according to the Method for Determining CF3:CF2.\nIn one embodiment, the article has a Quality Factor of at least about 0.25/mmH2O, (preferably at least about 0.5/mmH2O, more preferably at least about 1/mmH2O).\nIn some embodiments, the article includes a nonwoven polymeric fibrous web. Examples of suitable fibers for the nonwoven polymeric fibrous web include polycarbonate, polyolefin, polyester, halogenated polyvinyl, polystyrene, and combinations thereof. Particularly useful fibers include polypropylene, poly-(4-methyl-1-pentene), and combinations thereof. In one embodiment, the article includes meltblown microfibers.\nIn another aspect, the invention features an electret that includes a polymeric article having at least about 45 atomic % fluorine as detected by ESCA, and a CF3:CF2 ratio of at least about 0.45 as determined according to the Method for Determining CF3:CF2. In another embodiment, the electret includes at least about 50 atomic % fluorine as detected by ESCA, and a CF3:CF2 ratio of at least about 0.25 as determined according to the Method for Determining CF3:CF2.\nIn other aspects, the invention features a respirator that includes the above-described electrets. In still other aspects, the invention features a filter that includes the above-described electrets.\nIn one aspect, the invention features a method of making an electret that includes: (a) fluorinating a polymeric article to produce an article having surface fluorination; and (b) charging the fluorinated article in a manner sufficient to produce an electret. In one embodiment, the method includes charging the fluorinated article by contacting the fluorinated article with water in a manner sufficient to produce an electret, and drying the article. The method is useful for making the above-described electrets. In another embodiment, the method includes charging the fluorinated article by impinging jets of water or a stream of water droplets onto the fluorinated article at a pressure and for a period sufficient to produce an electret, and drying the article.\nIn other embodiments, the method includes fluorinating a polymeric article in the presence of an electrical discharge (e.g., an alternating current corona discharge at atmospheric pressure) to produce a fluorinated article. In one embodiment, the method includes fluorinating the polymeric article in an atmosphere that includes fluorine containing species selected from the group consisting of elemental fluorine, fluorocarbons, hydrofluorocarbons, fluorinated sulfur, fluorinated nitrogen and combinations thereof. Examples of suitable fluorine containing species include C5F12, C2F6, CF4, hexafluoropropylene, SF6, NF3, and combinations thereof.\nIn other embodiments, the method includes fluorinating the polymeric article in an atmosphere that includes elemental fluorine.\nIn other embodiments, the method of making the electret includes: (A) fluorinating a nonwoven polymeric fibrous web (i) in an atmosphere that includes fluorine containing species and an inert gas, and (ii) in the presence of an electrical discharge to produce a web having surface fluorination; and (B) charging the fluorinated web in a manner sufficient to produce an electret.\nIn other aspects, the invention features a method of filtering that includes passing an aerosol through the above-described electrets to remove contaminants.\nThe fluorinated electrets of the invention exhibit a relatively high oily mist resistance relative to non-fluorinated electrets.\nIn reference to the invention, these terms having the meanings set forth below:\nxe2x80x9celectretxe2x80x9d means a dielectric material exhibiting a quasi-permanent electrical charge. The term xe2x80x9cquasi-permanentxe2x80x9d means that the time constants characteristic for the decay of the charge are much longer than the time period over which the electret is used;\nxe2x80x9csurface modifiedxe2x80x9d means that the chemical structure at the surface has been altered from its original state.\nxe2x80x9csurface fluorinationxe2x80x9d means the presence of fluorine atoms on a surface (e.g., the surface of an article);\nxe2x80x9cfluorine containing speciesxe2x80x9d means molecules and moieties containing fluorine atoms including, e.g., fluorine atoms, elemental fluorine, and fluorine containing radicals;\nxe2x80x9cfluorinatingxe2x80x9d means placing fluorine atoms on the surface of an article by transferring fluorine containing species from a gaseous phase to the article by chemical reaction, sorption, condensation, or other suitable means;\nxe2x80x9caerosolxe2x80x9d means a gas that contains suspended particles in solid or liquid form; and\nxe2x80x9ccontaminantsxe2x80x9d means particles and/or other substances that generally may not be considered to be particles (e.g., organic vapors)."}
-{"text": "1. Field of the Invention\nThe invention relates to a readout device, and more particularly to an array-type readout device, a dual-function readout device, and a detecting circuit for a readout device.\n2. Description of the Related Art\nIn Wan-Jun Lin, Chao P. C. P., Shir-Kuan Lin, Hsiao-Wen Zan, \u201cA Novel Readout Circuit for an OTFD Gas Sensor with a New Front-end Trans-impedance Amplifier\u201d, Sensors, IEEE, pp. 1141-1144, 2011, an impedance detecting circuit is proposed. However, the proposed impedance detecting circuit includes two operational amplifiers, resulting in a large size that is unfavorable for use in an array-type readout device. Such a large detecting circuit is only suitable for use in a single-type readout device, and may have a relatively low sensitivity and a relatively low signal-to-noise (SNR) ratio."}
-{"text": "The field of the present invention relates to data mining techniques and, more particularly, to techniques for incorporating human interaction in an effective way so as to design similarity functions and perform class supervision of data.\nThe design of data mining applications has received much attention in recent years. Examples of such applications include similarity determination and classification. In the context of data mining, it is assumed that we are dealing with a data set containing N objects in a dimensionality of d. Thus, in this data space, each object X can be represented by the d coordinates (x(1), . . . x(d)). These d coordinates are also referred to as the features in the data. This is also referred to as the feature space which may reveal interesting characteristics of the data.\nThe effective design of distance functions used in similarity determination has been viewed as an important task in many data mining applications. The concept of similarity has been widely discussed in the data mining literature. A significant amount of research has been applied to similarity techniques such as, for example, those discussed in the literature: A. Hinneburg et al., xe2x80x9cWhat is the nearest neighbor in High Dimensional Space?,xe2x80x9d VLDB Conference, 2000; C. C. Aggarwal, xe2x80x9cRe-designing distance functions and distance based applications for high dimensional data,xe2x80x9d ACM SIGMOD Record, March 2001; and C. C. Aggarwal et al., xe2x80x9cReversing the dimensionality curse for similarity indexing in high dimensional space,xe2x80x9d ACM SIGKDD Conference, 2001, the disclosures of which are incorporated by reference herein.\nA different but related problem in data mining is the prediction of particular class labels from the feature attributes. In this problem, there is a set of features, and a special variable called the class variable. The class variable typically draws its value out of a discrete set of classes C(1), . . . C(k). A test instance is defined to be a data example for which only the feature variables are known, but the class variable is unknown. Training data is used in order to construct a model which relates the features in the training data to the class variable. This model can then be used in order to predict the class behavior of individual test instances, also referred to as class labeling. The problem of classification has been widely studied in the literature, e.g., J. Gehrke et al., xe2x80x9cBOAT: Optimistic Decision Tree Construction,xe2x80x9d ACM SIGMOD Conference Proceedings, pp. 169-180, 1999; J. Gehrke et al., xe2x80x9cRainForest: A Framework for Fast Decision Tree Construction of Large Data Sets,xe2x80x9d VLDB Conference Proceedings, 1998; R. Rastogi et al., xe2x80x9cPUBLIC: A Decision Tree Classifier that Integrates Building and Pruning,xe2x80x9d VLDB Conference, 1998; J. Shafer et al., xe2x80x9cSPRINT: A Scalable Parallel Classifier for Data Mining,xe2x80x9d VLDB Conference, 1996; and M. Mehta et al., xe2x80x9cSLIQ: A Fast Scalable Classifier for Data Mining,xe2x80x9d EDBT Conference, 1996, the disclosures of which are incorporated by reference herein.\nHowever, as sophisticated and, in some cases, complex as these similarity and classification techniques may be, these conventional automated techniques lack benefits that may be derived from human interaction during their design and application stages. Therefore, techniques are needed that effectively employ human interaction in order to design and/or perform data mining applications such as similarity determination and classification.\nThe present invention provides techniques for incorporating human or user interaction in accordance with the design and/or performance of data mining applications such as similarity determination and classification. Such user-centered techniques permit the mining of interesting characteristics of data in a data or feature space. For example, such interesting characteristics that may be determined in accordance with the user-centered mining techniques of the invention may include a determination of similarity among different data objects, as well as the determination of individual class labels. These techniques allow effective data mining applications to be performed in accordance with high dimensional data.\nIn accordance with a first aspect of the present invention, a computer-based technique of computing a similarity function from a data set of objects comprises the following steps/operations. First, a training set of objects is obtained. The user may preferably provide such training data. Next, the user is presented with one or more subsets of objects based on the training set of objects, wherein each subset comprises at least two objects of the data set. Preferably, the subset is a pair of objects from the data set. The user then provides feedback regarding similarity between the one or more subsets of objects. One or more sets of feature variables are defined based on features in the one or more subsets of objects. Next, one or more class variables are created in accordance with the user-provided feedback. Lastly, a similarity function or model is constructed which relates the one or more sets of feature variables to the one or more class variables.\nThus, advantageously, similarity between objects is represented as some function or algorithm determined by the attributes of the objects. The similarity model is then effectively estimated from the data set and user reactions.\nIn accordance with a second aspect of the present invention, a computer-based technique of classifying a test instance in accordance with a data set comprises the following steps. First, a test instance is obtained. The user may preferably provide such test instance. Next, the user is presented with at least one projection representing a distribution of the data set. The user then isolates a portion of the data presented in the at least one projection based on a relationship between the test instance and the data presented in the at least one projection. For instance, the user may isolate a subset of the data in the projection which the user determines to be most closely related to the test instance. Next, the behavior of the isolated portion of data is determined. Then, a class is determined for the test instance based on the isolated portion of data, when the user makes a decision to do so based on the determined behavior of the isolated portion of data. Alternatively, when the user makes a decision not to have a class determined for the test instance based on the isolated portion of data, other portions of the data set or a subset of the isolated portion of the data may be considered.\nFurther, in a preferred embodiment, the user is presented with two or more projections respectively representing different distributions of the data set such that the user may select one of the projections to be used when isolating a portion of data whose behavior is to be considered.\nThus, advantageously, such a class labeling methodology according to the invention provides a technique of decision path construction, in which the user is provided with the exploratory ability to construct a sequence of hierarchically chosen decision predicates. This technique provides a clear understanding of the classification characteristics of a given test instance. At a given node on the decision path, the user is provided with a visual or textual representation of the data in a small number of sub-spaces. This can be used in order to explore particular branches, backtrack or zoom-in into particular sub-space-specific data localities which are highly indicative of the behavior of that test instance. This process continues until the user is able to construct a path with successive zoom-ins which is sufficiently indicative of a particular class. The process of zooming-in is done with the use of visual aids, and can isolate data localities of arbitrary shapes in a given sub-space.\nIt is to be appreciated that the classification techniques of the present invention are more powerful than any of the conventional classification methods, since the invention uses a combination of computational power and human intuition so as to maximize user understanding of the classification without sacrificing discriminatory power. The result is a technique which, in most cases, can classify a test instance with a small amount of user exploration.\nThese and other objects, features and advantages of the present invention will become apparent from the following detailed description of illustrative embodiments thereof, which is to be read in connection with the accompanying drawings."}
-{"text": "(i) Technical Field\nThe present invention relates to an electronic conference assistance method and an information terminal device employed in an electronic conference system.\n(ii) Related Art\nConventionally, there is available an electronic conference system which has a large-scale readable and writable touch panel display device. Generally, such a touch panel display device is placed so as to be viewed by all participants of the conference and written thereon. Use of the electronic conference system enables a conference of a style, for example, where the participants gather around the touch panel display device, rather than remain seated, to discuss an idea conceived during the conference while writing the idea and so forth on the touch panel display device. The content written on the panel display can be stored intact as a screen image. Also in view of enhancement of conference efficiency, an increasing number of companies are introducing such electronic conference systems.\nMoreover, when such a system is employed, presence of all participants in the conference room where the touch panel display device is installed is not mandatory. That is, when another touch panel display device is installed in a conference room in another location and connected via a network to the electronic conference system in the main location, a remote conference can be realized. This allows a person in a remote location to participate in the conference.\nFurther, when a person who is supposed to participate in the conference but is away from the place where the conference is held as, for example, they are on a business trip connects their own personal computer (PC) to the electronic conference system, they can participate in the electronic conference from any desired place. Still further, when a portable phone is connected to the electronic conference system via a connection line, that person can participate in the conference through audio.\nAs described above, use of an electronic conference system can realize a conference of a style where participants can participate in a variety of manners, not limited to a conventional general conference style in which participants are kept seated and discuss ideas.\nHere, when a remote conference is taking place by connecting the device in the main location to the device used by a conference participant in another location to via a network, basically, data of images captured using a camera or the like in the conference room, in particular data of an image of a person who speaks, is transmitted to other locations. With this arrangement, a participant in the remote location can talk to the person speaking while looking at their image being shown.\nMoreover, a moderator who presides over the conference can check a participant in a conference by looking at the images captured in the respective locations when discussion is carried out between distant locations, and ask the checked participant to present their opinion or encourage them to speak."}
-{"text": "In an unassisted GPS-type position determination system, subscriber stations determine their own positions from satellite transmissions originating from the GPS-type position determination system, without requesting significant acquisition or calculation assistance from other network entities, for example, dedicated servers. That places significant processing demands on the subscriber stations because of the uncertainty in the timing, position, and frequency of these transmissions, requiring the subscriber stations to expend significant processing resources in searching for and locating these transmissions by, for example, testing large numbers of hypotheses varying the assumed timing, position and frequency of the transmissions. Since the number of hypotheses that must be tested is often staggering, the time required to search for the transmissions can be inordinately long and consume an excessive amount of processing resources, even for subscriber stations with dedicated receiver chains.\nThe uncertainty experienced by the subscriber stations stems from several sources. Assuming GPS positioning, there is first the uncertainty in knowing which of the 32 GPS satellites are visible to the subscriber station. That uncertainty is present because a subscriber station, upon power up or before a position fix is available, has no basis for identifying which signals of these 32 satellites can be usefully received. The useful reception of satellite signals is referred to as an ability of the subscriber station to \u201csee\u201d the satellite emitting the signal, or, in other contexts, as the satellite being \u201cvisible\u201d to the subscriber station.\nThis leads to inefficient searching because the subscriber station may waste considerable resources in searching for transmissions from satellites that are not visible to it, and which are therefore not useful for position determination purposes. For example, referring to FIG. 1, while satellites 54a, 54b, and 54c are visible to subscriber station 50 located at position 51 on the earth's surface 52, satellites 56a, 56b, and 56c are invisible to subscriber station 50, as they are located on the other side of the earth. Therefore, it would be wasteful for subscriber station 50 to search for the transmissions from satellites 56a, 56b, and 56c during a position fix attempt.\nIn addition, there is an uncertainty in knowing the timing or phase of the 32 chip PN \u201cgold\u201d codes that are embedded within the individual satellite transmissions. As these codes are circularly shifted versions of one another, the phase of a code uniquely identifies which of the satellites originated the transmission. The phase also reflects the propagation delay caused from transmission from the satellite to the subscriber station. To account for the possible variations in phase, the subscriber station must expend resources in searching over the full range of possible PN codes within a code phase searching window that is large enough to encompass the possible variations.\nMoreover, there is an uncertainty in knowing the relative movement between the subscriber station and the GPS satellites, which typically introduces a Doppler shift of approximately \u00b14 kHz in the frequency of transmission. To account for the possible variation of frequency introduced by the Doppler shift, the subscriber station must expend resources is searching over the full range of possible transmission frequencies within a frequency searching window that is large enough to encompass the possible variations caused by the Doppler shift.\nFinally, there is the uncertainty in knowing the degree to which the local oscillator (LO) of the subscriber station is out of tune with the GPS carrier frequency. Upon power-up, for example, it is not uncommon for the LO frequency to differ from the GPS carrier frequency by as much as \u00b15 ppm. Until synchronization between the LO frequency and GPS carrier frequency is achieved, the subscriber station must account for this uncertainty by increasing the size of the frequency search window that is employed.\nEven if the host wireless communications system or GPS-type position determination system eliminating some of this uncertainty by providing timing, positional information, or synchronization to the subscriber station, the processing demands on the subscriber station are often still substantial. For example, a synchronous system, such as a CDMA system, provides the subscriber station with time, and also synchronizes the LO frequency of the subscriber station to the GPS carrier frequency. Although the synchronization substantially reduces the LO frequency uncertainty, for example, from \u00b15 ppm to \u00b10.2 ppm, and the timing information allows the subscriber station to determine the position of the GPS satellites (using the GPS almanac or ephemeris data provided by the satellites), the subscriber station is still unable to determine which of the GPS satellites are visible to it, and it is still subject to the frequency uncertainty caused by Doppler shift."}
-{"text": "The present invention relates to an X-ray image intensifier.\nAs X-ray image intensifiers (to be referred to as \"I.I.\"s hereinafter), a general-purpose single visual field type I.I. and a high-grade variable visual field type I.I. are frequently used. In general, an I.I. comprises a vacuum housing which includes a substantially cylindrical outer casing, and an X-ray entrance window and an X-ray exit window which are arranged to close two ends of the outer casing. In the vacuum housing, input and output surfaces are arranged along the entrance and exit windows, respectively, and a focusing electrode constituting an electronic lens is located between the input and output surfaces. The I.I.s are classified into the single visual field type and variable visual field type due to differences in the number and arrangement of focusing electrodes, and the like. In the case of a variable visual field type I.I., when a voltage distribution to the focusing electrodes is switched, an output visual field image can be enlarged like, a normal visual field, a second visual field, a third visual field,....\nThe input surface has a base and a phosphor screen formed on the base, and has an arcuated circular shape.\nIn U.S. Pat. No. 3,716,713, the thickness of the phosphor screen is increased from its center toward the periphery, and is maximized at the periphery.\nAccording to an I.I. disclosed in Japanese Patent Disclosure No. 53-102663, the phosphor screen has the same arrangement as that in the above U.S. Pat. No., and the base has a mosaic structure having a large number of grooves for effecting a light guide function.\nAccording to an I.I. disclosed in Japanese Patent Disclosure No. 59-207551, the thickness of the phosphor screen is decreased from its center toward the periphery, and X-ray optical path lengths passing through the phosphor screen are adjusted to be equal to each other at the center and the periphery of the phosphor screen.\nIn the I.I.s having the above-mentioned arrangements of the input surfaces, the characteristic of an image obtained at the output surface, in particular, a luminance distribution characteristic, is such that a luminance is high at the center of the image and is gradually decreased toward the periphery. Therefore, a luminance distribution curve obtained as a result of measurement along the diameter of an image becomes a quadratic curve. In the variable visual field I.I., the same luminance distribution characteristic is obtained either in a normal visual field operation or in an enlarged visual field operation.\nThe reason for the above-mentioned luminance distribution can be considered as follows.\nIn the I.I.s disclosed in U.S. Pat. No. 3,716,713 and Japanese Patent Disclosure No. 53-102663, in order to prolong an X-ray passage distance in the phosphor screen, which influences light emission, so as to compensate for a quantity of light emitted from the phosphor screen, the thickness of the peripheral portion of the phosphor screen is increased. However, a portion between the intermediate portion and periphery of the phosphor screen cannot provide a similar effect upon increase in thickness, and, to the contrary, the luminance of the periphery of an image is decreased. This is because an excessive increase in thickness at the peripheral portion of the phosphor screen does not contribute to light emission of the phosphor by means of X-rays but degrades a transmittance of X-rays.\nIn Japanese Patent Disclosure No. 59-207551, in order to obtain a constant passage distance of X-rays at respective positions in the phosphor screen, the thickness of the phosphor screen is decreased at a given rate from its center toward the periphery. However, in order to obtain a theoretical luminance, the phosphor screen must be formed to have a uniform structure and a uniform emission intensity distribution. If these conditions cannot be satisfied, the luminance at the peripheral portion of an image, in particular, an area shifted from the center of the image toward the periphery by a distance 80 to 95% of an effective image diameter, is considerably decreased as compared to the above two prior arts.\nWhen the I.I.s having the above luminance distribution characteristic are used, the following problems are posed. In the distribution characteristic, the luminance at the center of an image is high and is decreased toward the periphery. When the I.I. is coupled to an optical system, a luminance difference between the center and the periphery of the image is emphasized due to an operation of the optical system. For this reason, a dark portion at the peripheral portion of the image has degraded discriminating ability of an object, and cannot be used for observing an object. Therefore, a virtual image area is decreased. When an object is observed upon clinical examination, a contour image of the object must be confirmed. However, when the effective image area is small as described above, the I.I. must be moved stepwise so that a portion to be observed is located at the center of the image. For this reason, the observation requires a long time, and an X-ray irradiation time is also prolonged. For example, when an observation is performed using a TV fluoroscopic imaging method, the entire object, i.e., the entire image, must be scanned, and this requires still more time.\nIn the enlarged visual field operation mode, e.g., in the second visual field operation mode, the luminance distribution characteristic of an output image is such that the center of an image is bright and the peripheral portion thereof is dark as in the normal visual field operation mode. In any visual field operation mode, an area of an input visual field is changed, but an image area which can be observed is almost not changed. For this reason, when the enlarged visual field operation is performed in order to microscopically observe the object after the contour image of the object is confirmed, the I.I. must be moved to locate the object at the center of image. If the object is a moving body, and is moved to the peripheral portion of an output image, the object cannot be discriminated since the luminance of the peripheral portion is low.\nSince the luminance distribution characteristic is not changed in the enlarged visual field operation mode, a low luminance portion is moved upon switching of visual fields. The object is often out of sight upon switching of the visual fields, and the I.I. must be moved to confirm the object at that time. For example, upon clinical examination wherein a change in object must be immediately judged, such as blood vessel imaging, the lack of necessary data and the complicated operations as described above may cause serious problems."}
-{"text": "The present invention relates generally to microlasers and associated fabrication methods and, more particularly, to Q-switched microlasers and associated fabrication methods.\nModern electro-optical applications are demanding relatively inexpensive, miniaturized lasers capable of producing a series of well-defined output pulses. As such, a variety of microlasers have been developed which include a microresonator and a pair of at least partially reflective mirrors disposed at opposite ends of the microresonator to define a resonant cavity therebetween. The microresonator of one advantageous microlaser includes an active gain medium and a saturable absorber that serves as a Q-switch. See, for example, U.S. Pat. No. 5,394,413 to John J. Zayhowski, which issued on Feb. 28, 1995, the contents of which are incorporated in their entirety herein. By appropriately pumping the active gain medium, such as with a laser diode, the microresonator will emit a series of pulses having a predetermined wavelength, pulse width and pulse energy.\nAs known to those skilled in the art, the wavelength of the signals emitted by a microlaser is dependent upon the materials from which the active gain medium and the saturable absorber are formed. In contrast, the pulse width of the laser pulses emitted by a conventional microlaser is proportional to the length of the resonator cavity. As such, longer resonator cavities will generally emit output pulses having greater pulse widths. Further, both the pulse energy and average power provided by a microlaser are proportional to the pulse width of the pulses output by the microlaser. All other factors being equal, the longer the microresonator cavity, the longer the pulse width and the greater the pulse energy and average power of the resulting laser pulses.\nConventional microlasers, such as those described by U.S. Pat. No. 5,394,413, are end pumped in a direction parallel to the longitudinal axis defined by the resonator cavity. In this regard, the longitudinal axis of the microresonator cavity extends lengthwise through the resonator cavity. Since the resonation cavity is generally a rectangular solid, the longitudinal axis is oriented so as to be orthogonal to the pair of at least partially reflective mirrors that define the opposed ends of the resonant cavity. As such, conventional microlasers are configured such that the pump source provides pump signals in a direction perpendicular to the at least partially reflective mirrors that define the opposed ends of the resonant cavity. The effective length of the resonator cavity is therefore equal to the physical length of the resonator cavity.\nWhile the microlaser can be fabricated such that the resonator cavity has different lengths, a number of factors contribute to generally limit the permissible length of the resonator cavity. See, for example, U.S. Pat. No. 5,394,413 that states that the resonator cavity, including both the saturable absorber and the gain medium, is preferably less than two millimeters in length. In particular, a number of electro-optical applications require microlasers that are extremely small. As such, increases in the length of the resonator cavity are strongly discouraged in these applications since any such increases in the length of the microresonator cavity would correspondingly increase the overall size of the microlaser.\nIn addition, the length of passively Q-switched microlasers is effectively limited by the requirement that the inversion density must exceed a predetermined threshold before lasing commences. As the physical length of the resonator cavity increases, greater amounts of pump energy are required in order to create the necessary inversion density for lasing. In addition to disadvantageously consuming more power to pump the microlaser, the increased pumping requirements create a number of other problems, such as the creation of substantially more heat within the microlaser which must be properly disposed of in order to permit continued operation of the microlaser. In certain instances, the heat generated within the microlaser may even exceed the thermal capacity of the heat sink or other heat removal device, thereby potentially causing a catastrophic failure of the microlaser.\nSince the pulse width and correspondingly the pulse energy and average power of the pulses output by a microlaser cavity are proportional to the length of the resonator cavity, the foregoing examples of practical limitations on the length of the resonator cavity also disadvantageously limit the pulse width and the corresponding pulse energy and average power of the pulses output by conventional microlasers. However, some modem electro-optical applications are beginning to require microlasers that emit pulses having greater pulse widths, such as pulse widths of greater than 1 nanosecond and, in some instances, up to 10 nanoseconds, as well as pulses that have greater pulse energy, such as between about 10 xcexcJ and about 100 xcexcJ, and greater average power, such as between 0.1 watts and 1 watt. As a result of the foregoing limitations on the length of the resonator cavity and the corresponding limitations on the pulse widths, pulse energy and average power of the pulses output by the conventional microlasers, conventional microlasers do not appear capable of meeting these increased demands.\nA microlaser is therefore provided according to one embodiment of the present invention that is capable of supporting a zig-zag resonation pattern in response to pumping of the active gain medium so as to effectively lengthen the microresonator cavity without having to physically lengthen the microresonator cavity. As such, the microlaser of these embodiments can generate pulses having greater pulse widths and correspondingly greater pulse energies and average power levels than the pulses provided by conventional microlasers of a similar size.\nAccording to the present invention, the microlaser includes a microresonator having an active gain medium and a Q-switch, such as a passive Q-switch proximate to and, in one embodiment, immediately adjacent to the active gain medium. In advantageous embodiments, the active gain medium and the Q-switch are integral such that the microresonator may be a monolithic structure. The microresonator extends lengthwise between opposed end faces. The microlaser also includes first and second reflective surfaces disposed proximate respective ones of the opposed end faces to define a microresonator cavity therebetween. While the first and second reflective surfaces can be coated upon respective ones of the opposed end faces of the microresonators, the first and second reflective surfaces can also be formed by mirrors that are spaced from respective ones of the opposed end faces. The microlaser can also include a pump source for introducing pump signals into the active gain medium via at least one of the end surfaces of the microresonator such that the zig-zag resonation pattern is established within the microresonator cavity.\nIn one advantageous embodiment, the opposed end faces are each disposed at a nonorthogonal angle xcex1, such as between about 30xc2x0 and about 45xc2x0, relative to a line perpendicular to a longitudinal axis defined by the microresonator cavity and extending between the opposed end faces. In one embodiment, the opposed end faces are each disposed at the same nonorthogonal angle xcex1 relative to the longitudinal axis such that the opposed end faces are parallel. In another embodiment, the opposed end faces are oriented in opposite directions by the same nonorthogonal angle xcex1. As a result of the nonorthogonal relationship of the opposed end faces, the microlaser of either embodiment is capable of supporting the zig-zag resonation pattern in response to pumping of the active gain medium via at least one of the end surfaces of the microresonator.\nBy supporting the zig-zag resonation pattern, the effective length of the microresonator cavity is increased relative to conventional microlasers having substantially the same physical size that do not support a zig-zag resonation path. In this regard, the effective length of the microresonator cavity of the present invention is the length of the zig-zag resonation path established by the microlaser which is significantly longer than the linear resonation paths established by conventional microlasers that extend parallel to the longitudinal axis of the resonator cavity. As such, the microlaser of the present invention can emit pulses having a longer pulse width and correspondingly greater pulse energies and average power levels than the pulses emitted by conventional microlasers of the same physical size.\nIn order to permit the pump signals to be received by the active gain medium without being reflected from the end face, the microlaser can include an antireflection coating on the end face through which the pump signals are delivered for permitting pump signals having a predetermined range of wavelengths to be received by the active gain medium. The microresonator also generally includes a plurality of side surfaces extending between the opposed end faces. In order to further facilitate resonation within the microresonator cavity, the plurality of side surfaces can be roughened, such as by grinding, to thereby diffuse light.\nIn order to permit the microlaser to emit signals of a predetermined lasing wavelength via one of the opposed end faces, the first reflective surface is preferably highly reflective for laser signals having the predetermined lasing wavelength. In contrast, the second reflective surface is preferably only partially reflective for laser signals having the predetermined lasing wavelength. As such, the microlaser can emit laser pulses having the predetermined lasing wavelength via the second reflective surface.\nIn one embodiment, the microlaser also includes a heat sink upon which at least the microresonator is mounted and a housing in which at least the microresonator is disposed. In this embodiment, the housing includes a window through which laser signals generated by the microresonator are emitted."}
-{"text": "An example of prior art chip capacitors is shown in FIG. 1. The chip capacitor shown in FIG. 1 includes a solid-state tantalum capacitor element 2 with a cathode layer 4 disposed on its outer surface. An anode lead 6 is led out from one end surface of the capacitor element 2. A flat cathode terminal 8 is connected to the cathode layer 4 with an electrically conductive adhesive (not shown). Also, a flat anode terminal 10 is welded to the tip end of the anode lead 6. An encapsulation 12 is provided by transfer molding with epoxy resin. Outer end portions of the flat anode and cathode terminals 10 and 8 are bent to extend along the end surfaces of the encapsulation 12 and, then, further bent to extend along the bottom surface of the encapsulation 12.\nIt is seen that a large proportion of the cathode terminal 8 is within the encapsulation 12, and the proportion of the volume occupied by the cathode terminal 8 to the entire volume of the encapsulation 12 is large. Further, both the cathode terminal 8 and the anode terminal 10 include portions extending on the side surfaces of the encapsulation 12. Accordingly, the length of the capacitor is increased by the thickness of these portions. In mounting such chip capacitor on a printed circuit board, the side surfaces of the cathode and anode terminals 8 and 10 are connected to the board by solder 14. Accordingly, when a number of such chip capacitors are to be mounted on a board side by side, as shown in FIG. 2, the spacing between adjacent chip capacitors must be large enough to prevent short-circuiting of adjacent capacitors, which prevents dense packing of the capacitors. Recently, smallsized, portable electric and electronic devices, such as cellular phones, have been remarkably improved, and chip capacitors to be used in such devices are required to be down-sized. For down-sizing prior art chip capacitors like the ones described above, the volume occupied by the capacitor element 2 in the chip capacitor including the encapsulation 12 should be as small as possible, which sometimes prevents the chip capacitor from having desired capacitance.\nTherefore an object of the present invention is to provide a chip capacitor which makes high density packing possible, and can have desired capacitance, while being small in size."}
-{"text": "A large and growing population of users is enjoying entertainment through the consumption of digital media items, such as music, movies, images, electronic books, and so on. The users employ various electronic devices to consume such media items. Among these electronic devices (referred to herein as user devices) are electronic book readers, cellular telephones, personal digital assistants (PDAs), portable media players, tablet computers, netbooks, laptops and the like. These electronic devices wirelessly communicate with a communications infrastructure to enable the consumption of the digital media items. In order to wirelessly communicate with other devices, these electronic devices include one or more antennas.\nThe conventional antenna usually has only one resonant mode in the lower frequency band and one resonant mode in the high-band. One resonant mode in the lower frequency band and one resonant mode in the high-band may be sufficient to cover the required frequency band in some scenarios, such as in 3G applications. 3G, or 3rd generation mobile telecommunication, is a generation of standards for mobile phones and mobile telecommunication services fulfilling the International Mobile Telecommunications-2000 (IMT-2000) specifications by the International Telecommunication Union."}
-{"text": "1. Field of the Invention\nThe invention relates to a drive unit, in particular for an injection unit or an ejector of an injection molding machine.\n2. Description of Related Art\nRecently, one has provided injection molding machines with electric and hydraulic drives, wherein actuations at high speed are exerted by the electric drive with relatively low forces, while the hydraulic drive is particularly advantageous if high axial forces have to be applied with comparatively minor actuations.\nIn the case of a closing unit of a plastics injection molding machine, for instance, the drive unit moves a movable tool faceplate of the machine. In so doing, the drive unit has to fulfill two important, different objects. On the one hand, it is to move the tool faceplate as quickly as possible for closing and for opening the mould so as to keep the cycle time of the manufacturing of an injection-molded component as short as possible. On the other hand, it is to impact the tool faceplate with a high clamping force, so that the tool can be kept shut against the high inner pressure during injection molding. The drive unit therefore has to be configured such that it is adapted to perform actuations at high speed and to apply high forces with a comparatively minor stroke. Requirements of this kind are posed, except with a closing unit, also with the actuation of ejectors or the injection unit of an injection molding machine.\nDE 101 21 024 A1 (cf. in particular FIGS. 26, 34) of the Applicant discloses a drive unit that is adapted to fulfill the afore-mentioned requirements. This drive unit comprises a hydraulic force transmitting element, the smaller piston unit of which is actuated via an electrically actuated stroke spindle device for closing a tool. This smaller piston unit may consist of one single smaller piston, or of a plurality of small pistons. These confine, along with a cylinder or interface and one or several large pistons of the force transmitting element, a pressure chamber, wherein, by the moving of the small piston unit into the pressure chamber, a high pressure can be generated, which acts, via the large active surface of the large pistons (power pistons) on the movable tool faceplate which may then be kept shut with high force. During the quick closing of the tool with comparatively low force, the interface is indirectly connected with a spindle nut of the spindle device, so that the piston unit with smaller diameter, the power piston, and the interface are jointly shifted by the spindle device. For applying the high force, the interface is fixed at the frame of the injection molding machine, so that the further closing movement of the tool is determined by the moving of the smaller piston unit into the pressure chamber and the corresponding axial movement of the large piston of the force transmitting element.\nIn one embodiment described in DE 101 21 024 A1 (FIG. 34), the coupling of the cylinder to the stroke spindle device is performed hydraulically. To this end, a chamber confined by a section of the small piston unit and the cylinder is impacted with pressure from a high pressure storage means, so that the pressure medium incorporated in the chamber acts like rigid pulling mechanics and the cylinder participates in the closing stroke of the stroke spindle device and thus of the small piston unit.\nIn an embodiment illustrated in FIG. 26 of DE 101 21 024 A1, the small piston unit is, during rapid motion, connected with the large piston via an electromagnetic coupling. This large piston is in turn centered with respect to the cylinder by a prestressed centering spring arrangement. The prestressing of this centering spring arrangement is chosen such that the axial shifting of the small piston unit is, during rapid motion, transferred to the large piston via the coupling, and from there via the centering spring arrangement to the cylinder so as to take it along.\nIn both known solutions the force transmitting element is designed to be double-acting, so that, for tearing open the tool, a high tear-open force acts on the tool via the force transmitting element as the small piston unit moves in opening direction. This movement of the small piston unit in opening direction is performed during the application of the tear-open force against the force of a prestressed pressure spring.\nA disadvantage of the initially mentioned known construction (FIG. 34) is that, for applying the high pressure in the chamber during rapid motion, a comparatively complex circuitry with high pressure storage means and electrically controlled direction control valve is required, so that this circuitry variant is very expensive and also requires substantial construction space.\nIn the solution illustrated in FIG. 26 of DE 101 21 024 A1, the large piston has to be designed with a very large surface due to the integrated coupling, so that a compact solution cannot be realized with such a construction."}
-{"text": "1. Field of the Invention\nThe present invention relates to a communication device incorporating the MAP (Manufacturing Automation Protocol), an international standard communication protocol that has been defined in ISO Standard ISO/DIS 9506-1, which is useable in a factory automation (FA) environment.\n2. Description of the Background Art\nIn an automated factory, a variety of devices are employed in the manufacturing operation and the devices are joined through a local communication network into a factory system. Since certain devices may be more suitable than others to perform desired manufacturing operations, often the devices used in the factory system will be manufactured by different vendors. Accordingly, each such FA device, whether a factory computer, robot, numerical control (NC) machine, programmable logic controller (PLC), process control equipment, or the like, will have a different type of microprocessor, use different computer languages and execute customized programs. It is desirable that the internal processing and operation of each device should have little effect on the way the devices interact in the factory system and, in particular, how they communicate with each other. In order to provide a common basis for communication, all of the devices in the system must use a common message structure (\"syntax\") and use a common set of messages or \"semantics\" (i.e., the naming of and access to remote variables, program loading, job management, error reporting and the like).\nThe Manufacturing Message Specification (MMS) has been adopted as an international standard that permits programs to be written for a variety of factory system devices on the basis of common semantics and syntax. The MMS is specified in two parts comprising the message services (semantics) and the protocol (syntax). The message services are grouped into functional units that relate to the kinds of functions that are performed when an application (a program that performs some desired job) at one user location interacts with the local communication network for purposes of communicating with a user at another (remote) location. A total of 86 message services may be grouped according to the functions of context management (e.g., Initiate, Conclude, Abort, Reject, Cancel), remote variable services (e.g., Read-data, Write, Define Named Variable, etc.), program services (Initiate Download Sequence for a program, Load Domain Content, etc.), diagnostics (Status, etc.), operator communication (Input and Output), coordination between applications (Define Semaphore, etc.), file services (File Open, File Read, etc.), event management (Define Event Condition, etc.), journal management (Read Journal, Write Journal, etc.) and job management/device control (e.g., Start-robot movement, Stop, Resume, etc.). A detailed description of the MMS standard appears in \"MMS Tutorial by John R. Tomlinson, System Integration Specialists Company, Inc. (1987).\nFIG. 4A is a block diagram illustrating the connection of two stations, each having corresponding applications and being interconnected by a local communication network, as they would appear in an automated factory environment. The application in station A \"at one end\" of the network communicates, via a MMS provider (shown as MMS), a logic link controller (LLC), a media access controller (MAC) and a modem at each station that is connected to a local network, with the application in station B \"at the other end\" of the network. In conventional MMS terminology, for such communication, station A is the \"Client\" and requests station B as the \"Server\" to perform some application specific operation; the Server responds with information resulting from the operation as it is performed. Typically, the Client is a controller station and the Server is a FA device.\nFIG. 4B is a block diagram illustrating the arrangement of a conventional communication device employing a PLC (Programmable Logic Controller) 1 as an example of an FA (Factory Automation) device. Ordinarily, the PLC has limited storage capability and relies on outside storage media (e.g., disk storage) to store pertinent programming and variables. Reliance on outside storage media has the disadvantage that when a power failure or OFF condition is encountered by the PLC, the relationship between the PLC and its external storage media must be redefined at power ON.\nIn FIG. 4B, the numeral 2 indicates a MAP interface unit, serving as a communication device and being connected between a MAP network 3 and the PLC 1 via a PLC-dedicated bus 4. The MAP interface 2 comprises an MMS protocol 5 whose communication object is a named variable, rather than an address. A PLC driver 7 for accessing the PLC 1, and a local manager 8 for carrying out management functions also are found in the MAP interface unit 2.\nFinally, interface 2 includes a VMD (Virtual Manufacturing Device) 6 for converting the MMS protocol 5 into a protocol reflecting the resources and functionality of the real FA device, e.g., PLC 1 in the preferred embodiment, and performing a process corresponding to each MMS service. The VMD, as an abstract representation of a Server showing its external behavior, comprises four conventional abstract elements including Executive function, Capabilities, Program Invocations and Domains. The latter are dynamic in nature and come into existence and are removed from the system either by MMS Services or by local action. The Domains comprise instructions and/or data which is dedicated to specific resources, such as the portion of the machine or robot that is controlled. Services are provided for a Client to manipulate Domains that are defined at the MMS Server, such as the Initiate Download Sequence and Upload Segment services.\nIn the standard MMS specification, the Domain management services comprise a Domain Object attribute, which specifies a VMD Object-specific name or Domain Name to uniquely identify the Domain within the VMD, and a List Of Capability attribute, which is a list of implementation specific parameters necessary to partition the resources of the VMD.\nThe PLC 1 is equipped with a computer interface 11. PLC 1 includes a symbolic address variable registration section 12, and is connected to the MAP network 3 via the dedicated bus 4 and MAP interface unit 2. A controller and multiple FA devices, each representing a different station, may be connected to the MAP network 3 for communication therebetween.\nThe VMD 6, as a \"virtual device\" that serves as an abstract model of the MMS server application, provides a consistent basis for defining the MMS services for all devices. In the present case, VMD 6 models the externally visible behavior of the PLC 1 and comprises applications that provide several MMS services and are represented as units, including Define Named Variable/Delete Named Variable means 61 for defining and deleting a named variable convertibly into a symbolic address variable specific to the PLC 1. Also included in VMD 6 is named variable accessing means 62 and a variable conversion table 63 wherein a named variable is registered (stored) in correspondence with a symbolic address variable specific to the FA device.\nFIG. 5 is a flowchart showing the operation of the MAP interface unit 2 acting as the communication device known in the art. The operation of the MAP interface unit 2 will now be described in reference to FIG. 5.\nReferring to FIG. 5, when a request for a Define Named Variable service is received from a station B at the other end (not shown) that is connected to the MAP network 3 in Step 201, the MMS protocol 5 activates the Define Named Variable/Delete Named Variable means 61 in the VMD 6 in Step 202. As a result, for example, the Define Named Variable/Delete Named Variable means 61 may register a named variable, e.g., \"DATA001,\" into the variable conversion table 63 in correspondence with a symbolic address \"D1\" according to the request of the other-end station B in Step 203. The named variable is related to a particular FA device, e.g., robot 1, as contrasted to robot 2 which may be represented by named variable \"DATA002\", and each FA device may be made by any of several vendors. Accordingly, the named variable is identified as having a relationship to a symbolic address, which ordinarily is vendor specific, e.g., Mitsubishi Electric Company of Japan has the standard address D1 and other unique standard addresses are assigned to other vendors. If the request is for a Delete Named Variable service, a corresponding named variable is deleted from the variable conversion table 63.\nWhen a request for a variable access service to the named variable \"DATA001\" is then received from the other-end station B in Step 204, the MMS protocol 5 activates the named variable accessing means 62 in the VMD 6 in Step 205, the named variable accessing means 62 converts the named variable \"DATA001\" into the symbolic address \"D1\" using the variable conversion table 63, and the VMD 6 accesses the symbolic address \"D1\" of the PLC 1 via the PLC driver 7 in Step 206. By using the table 63 which defines a named variable (e.g., DATA001) to be a vendor specific address (e.g., D1) programming is simplified and is useable for any of several devices from different vendors, since only a data call that is generic to the FA device at a given location (i.e., using DATA001) is used in the program to identify a desired operation, rather than a particular vendor address.\nSince the named variable is defined in a procedure as shown in Steps 201 to 203, i.e., when a request for the Define Named Variable service is received from the other-end station B (not shown), the MMS protocol 5 activates the Define Named Variable/Delete Named Variable means 61 in the VMD 6 to cause the Define Named Variable/Delete Named Variable means 61 to register the named variable into the variable conversion table 63 in response to the request of the other-end station B, registration cannot be made from other than the other-end station B. Accordingly, an application concerning a named variable to be registered for the other-end station B must be added for registration.\nSeveral other problems also are encountered in the conventional system design. For example, while a total of 86 services are set forth in the MMS, services which historically experience a low request level may not be provided. In fact, the actually provided services often comprise only about half of the total available services, due to the limited memory capacity in the PLC. For example, the other-end station B often is not provided with the Define Named Variable service or with the Delete Named Variable service. In the absence of these services, the table 63 cannot be utilized effectively, particularly when a power outage or OFF condition is encountered.\nMoreover, the known communication device arranged as described above does not allow a user-defined named variable for accessing an FA device to be registered from other than the other-end station. This requires an application for registering the named variable to be added to the other-end station for the purpose of registration."}
-{"text": "Frequency division multiplexing enables the concurrent communication of multiple signals over the same physical medium. In a frequency division multiplexed system, signals are frequency-converted to an assigned frequency band prior to being transmitted over the physical medium. To enable recovering the signals at the receiver, each of the different signals is assigned to a different frequency band or bands. The receiver then separates the received composite signal into the various frequency bands, and then processes the signal received in one or more of the assigned frequency bands to recover the information contained in that signal. Conventional circuitry utilized for separating the frequency bands, however, is costly.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with some aspects of the present invention as set forth in the remainder of the present application with reference to the drawings."}
-{"text": "The present disclosure relates to a vehicle, and more particularly, to a personalized route planning system therefore.\nVehicles often include computer-implemented mapping systems. The mapping systems typically include route planning applications to provide users with directions between different locations. The route planning application includes representations of roads and intersections and one or more algorithms to output a suggested route of travel. These algorithms can output routes depending upon user-selected parameters. For instance, a route planning application can enable a user to select a time efficient route, or a distance efficient route.\nOver the last several years, users have grown to rely increasingly on route planning applications. Personalized tailoring of such routes, however, has been deficient."}
-{"text": "The present invention relates to the field of folding doors with flexible door leaves. More specifically, the invention relates to a door comprising a door leaf which is at least partly made of a flexible cloth material and which is movable between a closed position and an open, folded position, in which the door leaf is folded around a plurality of folding lines extended between opposite side edges of the door leaf, a plurality of guide members which are connected to the opposite side edges in a spaced-apart relationship along the same, and two side frames which extend adjacent to a respective side edge for guiding the guide members. Such a door is known from e.g. EP 0 113 634. The invention also relates to a method for assembling such a door.\nSince the 1970s there has been a great need to use rapidly moving doors in buildings for industrial use. This applies to openings indoors as well as in external walls, where the door provides shielding between different activities or prevents draughts/heat losses. Presently, rolling doors with flexible door leaves are used for this purpose, which doors are rolled up on an overhead drive shaft and which can be provided with transverse wind reinforcements on the door leaf to counteract wind load. For security reasons, rolling doors can be provided with a safety edge protection, a drop protection, etc.\nAlongside the development of rolling doors, there has been a development in foldable doors according to the introductory paragraph, in which the door leaf is instead folded as it is lifted during the opening process. These door leaves, too, are often provided with transverse wind reinforcements, comprising beams or sections which are suitably connected to the flexible door leaf. The wind reinforcements also contribute to the lateral stability of the door leaf.\nThe lifting arrangements of known folding doors vary from case to case, but usually the door leaf is lifted with the aid of at least one pair of belts/wires in the lowermost section, so that the transverse sections are gradually gathered in a bundle when the door is opened.\nEP 0 113 634 describes a folding door with transverse reinforcement sections. Every other section, beginning with the lowermost one, is extended into the side frames and supports guide rollers which are guided by the side frames in the depth direction, i.e. perpendicular to the door opening. The intermediate sections are shorter and have no guide rollers. Three lifting belts, which run vertically along the door leaf, are each connected to the bottom section. When the belts are rolled up on a transverse overhead shaft, they pull the bottom section upwards, which in turn successively pulls the other sections upwards so that the door leaf is folded in horizontal folds. Since every other section lacks guide rollers and consequently is not guided by the side frames, in the open position these non-guided sections will hang like a cradle by the intermediary of two superjacent guided sections, so that the door leaf is folded like a concertina. By virtue of the fact that the belts run on the exterior of the door leaf and on one and the same side thereof, all the non-guided sections are forced to fall out on the opposite side of the door leaf during the opening motion. Thus, in this known door, the lifting belts ensure that the non-guided sections fall out in one and the same direction.\nFR-A1-2,706,941 describes a folding door which, in conformity with the door in EP 0 113 634, has transverse reinforcement sections of which only every other section is guided by the side frames, and where the intermediate sections are non-guided in order to fall out sideways when the door is being closed. However, edge guide members are lacking, and the two side edges of the door leaf hang essentially completely unguided in the depth direction, received in the side frames. In this door, too, the lifting belts are used to ensure that the non-guided sections fall out on one and the same side of the door leaf. The lifting belts are located adjacent to the side frames.\nFR-A1-2,722,531 describes a door in which all the transverse reinforcement sections run in one and the same relatively wide guide track in the side frames and where the lifting belts are attached to the second lowest section and run through special belt loops in every other section. These loops result in the sections with loops gathering in a first bundle during lifting, while the sections without loops gather in a second bundle, hanging from the first bundle. The loops ensure that the sections without loops fall out on one and the same side of the door leaf in connection with lifting. Extra safety belts begin operating if the regular belts should break. All belts are located in the door opening between the side frames.\nSE 454,526 describes a technique for achieving forced folding of a door leaf, which is divided into horizontal, mutually foldable sections. In an embodiment shown in that document, the door leaf is designed in the form of a unitary, flexible piece of cloth, where every other section beginning with the lowermost is stiffened at its vertical side edges by means of rigid side borders. Every such rigid side border is provided with an upper and a lower guide pulley, which guide pulleys have a constant vertical relative position. These two pulleys run in an associated vertical guide track formed in the stationary side frame of the door. The guide track opening facing the door opening is provided with flanged edges for retaining the guide pulleys in the guide tracks. Thus, there is a plurality of guide tracks in each side frame. The number of guide tracks in each frame equals the number of sections provided with rigid side edges. Thus, only two guide pulleys run in each guide track, and, as a result of the stiffening, the stiffened sections are always vertically orientated in line with their associated guide track, and no folding takes place of these sections in connection with lifting. More specifically, the stiffened sections function as essentially completely rigid sections. In one example, the door leaf has three stiffened and three non-stiffened door leaf sections; and consequently three parallel guide tracks in each side frame.\nIn SE 454,526 mentioned above, two wires or the like are fastened to the lowermost, stiffened section for lifting and folding the door leaf. During lifting, the non-stiffened sections will be folded in between the stiffened sections, which assume a position beside each other like books on a shelf. When the door leaf has been lifted completely, a concertina-like bundle is obtained where the vertical, stiffened sections stand next to each other in a respective guide track and each intermediate, flexible section is extended obliquely downwards from the top of a stiffened section to the bottom of an adjacent stiffened section. In the lifted position, the whole bundle hangs from the section to which the wires are fastened.\nKnown folding doors of the type mentioned above exhibit various drawbacks depending upon the design chosen.\nIn the cases where the lifting belts and any associated loops are placed on the door leaf itself, there is a risk that individuals and vehicles will get caught in and lifted with the door leaf during opening. Moreover, such a placement is not aesthetically pleasing. Making holes for the lifting loops results in indication of fracture/weakening of the door leaf and additional manufacturing costs. In addition, centrally located lifting belts require a horizontal drive shaft or the like above the door.\nAnother drawback of the prior art doors is that the folding of the door leaf is effected in a non-reliable manner, or in a manner resulting in undesired wear of the door leaf. For example, the door leaf can be folded either inwards or outwards depending on the current pressure difference. This may, for example, result in the door leaf wearing against the upper edge of the door opening and/or the belts.\nAny pressure differences are absorbed by the transverse reinforcement sections, which, consequently, are squeezed against the side frames. In that way, in some known doors, the side edges of the door leaf are squeezed between the sections and the frames, resulting in the door leaf wearing out.\nMost known folding doors of the type described by way of introduction have a relatively wide side frame in the depth direction (i.e. transversely of the door opening) for receiving the side edge of the door leaf. Such a wide side frame is required to prevent the door leaf from jamming in the side frame during opening and closing. One drawback of having a wide side frame is that the door leaf can move in the depth direction in an undesired manner in connection with pressure differences, resulting in an undesired ability to move in the depth direction in the closed position, a poor aesthetic impression, and incomplete sealing. Moreover, a wide frame requires a large installation area and is expensive and heavy to make and assemble. A particular drawback of the door according to SE 454 526, wherein each stiffened section runs in its own guide track, is precisely that the side edges become very wide and costly as the height of the door and the number of sections increase, since a separate guide track is required for every other section of the door leaf.\nThese and other drawbacks of the prior art will appear clearly below in connection with the description of the invention.\nIn order to reduce the above-mentioned drawbacks of the prior art, according to the invention a door is provided of the type stated by way of introduction, i.e. a door comprising a door leaf which is at least partly made of a flexible cloth material and which is movable between a closed position and an open, folded position, in which open position the door leaf is folded about a plurality of folding lines extended between opposite side edges of the door leaf, a plurality of guide members which are connected to the opposite side edges in a spaced-apart relationship along the same; and two side frames which are extended adjacent to a respective side edge for guiding the guide members. The door according to the invention is characterised in that each side frame defines at least a first and a second guide groove, that said guide members comprise, at each side frame, a first set of guide members running in the first groove only of the side frame, and a second set of guide members running in the second groove only of the side frame, and that the first and the second guide members are connected to the door leaf in such a way that the side edges, in the folded position of the door leaf, run back and forth between the first and the second guide groove with said folding lines defined by the guide members.\nA xe2x80x9cflexible cloth materialxe2x80x9d could be any suitable kind of cloth, fabric or sheet of a flexible, foldable material, which can be coated or uncoated.\nWhen the door according to the invention is being opened or closed, the first guide members run in the first groove only and the second guide members run in the second groove only. In each groove, the associated guide members will be successively brought together during the opening motion. As a result, the mutual distance between the first guide members as well as the mutual distance between the second guide members will decrease when the door opens. Although, at present, it is probably preferable to have two guide grooves only in each side frame, it is within the scope of the invention to add one or more supplementary guide tracks, but in such variants it is still the case that the guide members in the first and the second guide groove are mutually brought together during the opening motion.\nThe expression xe2x80x9cguide groovexe2x80x9d can refer to a physical channel or the like, but it can also be interpreted as an abstract term and shall be considered to include all variants where the side edges are provided with special guide devices or means for defining two separate, predetermined movement paths or tracks for the guide members. Usually, the two guide grooves, which are defined by the side frames, are juxtaposed transversely of the door opening, but it is also possible that this distribution in the side frame itself is in a direction parallel to the door opening. In the latter case, there must be special connection members between the guide members and the edges of the door leaf, so that the attachment points in the edges of the door leaf run along two parallel lines or paths spaced from each other transversely of the door opening. In one embodiment, the first and the second guide groove can, for example, each be formed as a physical channel, whose side walls achieve the guiding of the guide members. These channels can be open towards the door opening but, with suitable connection members between the guide members and the door leaf edges, it is possible to turn the openings of the channels away from each other, so that one opening faces the front of the door and the other opening faces the rear of the door. As an alternative to physical channels, each guide groove can instead be defined by a rod or the like fixedly arranged in the side frame with which the guide members engage slidably in a suitable manner.\nUsually, the door according to the invention would be orientated with vertical side frames and a vertically guided door leaf. However, it is within the scope of the invention to place the door horizontally instead, but to facilitate the description and definition of the invention, terms such as xe2x80x9cliftingxe2x80x9d, xe2x80x9cvertical side framesxe2x80x9d, etc. are used throughout this specification. Accordingly, if the door is to be placed lying down, these orientation-determining expressions should be interpreted to include the horizontal case as well.\nIt should be noted that the above-mentioned xe2x80x9cplurality of guide membersxe2x80x9d can comprise xe2x80x9cfurther guide membersxe2x80x9d in addition to said first guide members and said second guide members, for example special guide members at the closing edge of the door leaf. Even if the first and the second guide members are normally located alternatingly in the first and the second guide groove, there may be portions of the door leaf where two adjacent guide members are located in the same guide groove.\nSeveral advantages are achieved by the invention by the provision of the double guide tracks in the side frames, as well as by the distribution of the guide members in the same:\n1. A first advantage of double guide tracks is that the folding of the door leaf becomes much more exact and controlled in comparison with how the folding takes place in the known doors. A controlled folding in the side frames in turn leads to generally safer functioning with a reduced risk of a breakdown, and to a considerable improvement in the appearance of the door leaf during operation. Moreover, no special pre-folding members or wear protection is necessary.\n2. As mentioned above, the side frames of the prior art doors must often be wide in the depth direction of the door opening in order to prevent the door leaf from jamming during lifting. The door according to the invention does not have that problem. Accordingly, a second advantage of double guide tracks is that the depth of the frame can be reduced considerably. The depth of the frame is mainly determined by the size of the guide elements in the depth direction of the door, but also by the amount of space required for the side edge itself of the door leaf.\n3. A third advantage of double guide tracks and of the side frames actively influencing the folding in the direction desired is that all lifting members, such as belts or wires, can be located protected within the side frames. Unlike in known doors, the lifting members need not be mounted on the surface of the door leaf for guiding the folding, but can be located protected in the side frames. This in turn means that both the lifting members and the environment are protected. The general appearance of the door also becomes more attractive with concealed lifting members. The driving can be achieved with two lifting points only, and if a variant with a transverse drive axle is used, it can be made with a less substantial dimension. Placing all the lifting members in the side frames also yields the advantage that no transverse drive shaft is needed above the door since belt drums can be attached directly to the side frames. However, it should be noted that, for example, in connection with very wide and/or heavy doors, it might be necessary to provide supplementary safety belts/lifting belts in the middle to prevent deflection. However, unlike in the prior art, it is not necessary to use such an additional belt for guiding the folding, but only for reducing the stress on any fall-out-preventing means in the side frame.\n4. A fourth advantage of double guide tracks is that, in its closed position, the door leaf can be positioned centrally in the depth direction between the guide tracks. This results in improved sealing and appearance, reduces wear and provides a more compact frame. In particular, the side edges of the door leaf can be guided in separate sections in the depth direction for obtaining an exceedingly compact door in the depth direction.\nNormally, the transverse folding lines, or extensions thereof, of the door leaf, will intersect the guide grooves. Accordingly, if the door leaf is provided with a plurality of transverse reinforcement members, each of which is extended between an associated pair of guide members, extensions of these reinforcement members can intersect the guide grooves for defining the folding lines of the door leaf. In order to obtain a straighter door leaf in the closed position, all the reinforcement members, or at least the majority of them, can lie alternatingly on the one and on the other side of the door leaf. In a special case, the two lowermost reinforcement members can be located in the same guide groove.\nWith respect to the space requirement at the upper part of the door, it will be appreciated that, in principle, the space required in the depth direction for the reinforcement members, when these are piled on top of each other according to the invention in two guide grooves, is only half as large as in the prior art where they are piled up in one and the same channel.\nIn a preferred embodiment of the invention, the first guide groove and the second guide groove in each side frame comprise a first physical guide channel and a second physical guide channel respectively, which are open in the direction of the door opening and have a width in the depth direction which is adapted to the dimensions of the guide members in the same direction. In this connection, the guide members can consist of non-rotatable sliding members or rollers. However, it is essential that no large play is required in the depth direction between the guide members and the side walls of the guide channels.\nEach guide channel can be provided with a fall-out-preventing means for retaining the guide members. The guide members can be designed themselves to prevent a fall if a lifting member breaks.\nAccording to a first embodiment, each side frame is provided with a U-section, whose bottom wall partly covers the two guide channels in order to form fall-out-preventing means, extended along the frame and open towards the door opening. In this connection, this U-section can have a double function since the side edge of the door leaf can be inserted in and seal against the U-section. One part of this U-section can be detachable for installation and maintenance. This embodiment yields the advantage that both the guide members and the side edges of the door leaf have a limited ability to move in the depth directionxe2x80x94i.e. they have good guiding in the depth direction and that the door leaf is centred in the depth direction relative to the guide tracks.\nAccording to a second embodiment, each side frame comprises a bottom wall, a first outer side wall and a first partition which both extend from the bottom wall for defining said first guide channel, and a second outer side wall and a second partition which both extend from the bottom wall for defining said second guide channel, wherein said first and second partitions define therebetween a space which receives the side edge of the door leaf. In this embodiment, the partitions serve two purposes: they define the guide channels and they receive therebetween the side edge of the door leaf in order to guide it along the side frame. In this embodiment, said partitions and said outer side walls can be provided with fall-out-preventing flanges adapted to retain the guide members in the guide channels.\nThe door leaf can be formed optionally as a continuous piece or divided into sections held together with e.g. transverse reinforcement sections. The door leaf can be formed entirely of a flexible cloth material, but the invention will also work if some door leaf sections are rigid. More specifically, the door leaf can be lifted in such a way that every other section is not folded, and these section can be made of a more rigid or a completely rigid material, while the other sections which are folded must be made of a flexible material.\nPreferably, there is at least a first flexible pulling member, such as a belt, a wire, a chain or the like, in each side frame for guiding the movement of the door leaf. If the guide grooves are physical channels, the pulling members can suitably be located in the same. In one embodiment, a direct lifting force is applied to only a single guide member in each side frame, called a driven guide member. If the door leaf runs alternatingly between the two guide tracks all the way down to its closing edge, the lifting can be effected in the lowermost guide member. However, in some cases, there may be a special safety arrangement with a bottom section having a reduced weight. In such a case, the lifting be effected in the second lowest guide member instead. If, however, there is only one pulling member in each side frame, these members can consist of a continuous pulling member. Moreover, there can be double pulling members or more in each side frame.\nIn principle, the lifting force applied to the driven guide members can be transmitted to superjacent guide members in two different ways. Either the design is such that the guide members strike against each other during the lifting, so that the lifting force is transmitted directly in the side frame. Alternatively, transverse reinforcement sections are used which are of such thickness that they will strike against each other before the guide members strike against each other. In this case, the lifting forces are instead transmitted by the intermediary of the reinforcement sections and, specifically, in this connection, guide members in the form of rotatable rollers can be used, which may be problematic if the guide members are to abut against each other.\nIn one embodiment of the invention there may also be a further pulling member in each side frame which applies a direct lifting force to a second driven guide member, the first and the second driven guide members running in different grooves. If, for example, the second driven guide member is located closer to the closing edge of the door leaf, its pulling member can be driven a somewhat longer distance than the first pulling member for achieving an xe2x80x9cextra liftxe2x80x9d of the second driven pulling member during the opening motion of the door. This can, for example, be achieved by the use of larger diameters in the winding drums for the second pulling members and/or greater thickness in the latter. Another possibility is to lift the last section more at the end of the lifting motion by virtue of only the lower part of the pulling member having a substantially greater thickness or to mount a member on the lower part of the pulling member which gives it an extra lifting motion at the end of the opening motion. The advantage of such an extra lift is that the vertical dimensions of the door leaf in the open position can be further reduced.\nFor easy transportation and installation of the door, each side frame can be divided into a shorter top part and a longer bottom part. The top parts are made with such a length that all guide members, which are connected to the side edges of the door leaf, can be received in the top parts simultaneously. In this way, the whole door leaf, all the reinforcement sections, all the lifting members, the upper part of the frames as well as the drive unit can be pre-assembled at the factory and be delivered to the installation site as a single unit. The top parts with the guide members inserted therein are mounted to the bottom parts only at the location where the door is to be installed. In assembling the frame parts, the guide grooves are likewise assembled, and the guide members and the pulling member can then be inserted into the side frames and the door can be used directly.\nThese and other embodiments and advantages of the invention will appear from the claims and from the following detailed description of preferred embodiments."}
-{"text": "1. Field of the Invention\nThe present invention generally relates to a method for directly writing data into an optic disk without a computer system; and in particular to a method for directly writing data retrieved from an electronic data storage into an optic disk without a computer system interfacing therebetween.\n2. The Related Art\nWith the rapid development and prevalence of electronic storages, such as compact flash memory, more and more data are stored in the electronic storage for portability and fast access. To more space-efficiently store the data, some people prefer to transfer the data from the electronic storage to optic disks for data backup purposes. Heretofore, the optic disk drive must be connected to a computer system for data writing operation. Thus, the data have to be read into the computer system and then written by the computer system into the optic disk accessed by means of the optic disk drive. This causes problems. For example, a computer system is a must in transferring data from a portable compact flash memory device to an optic disk. Thus, such a data transfer operation cannot be carried out without a computer system having proper data ports.\nThus, the present invention is aimed to solve the above problem by providing a method for directly transferring data from a compact flash memory device to an optic disk without a computer system interfacing therebetween.\nAccordingly, an object of the present invention is to provide a method for writing data retrieved from an electronic storage to an optic disk without a computer system interfacing therebetween.\nAnother object of the present invention is to provide an optic disk drive capable to perform a direct writing operation to an optic disk without being controlled by a computer system.\nTo achieve the above objects, in accordance with the present invention, there is provided a method for directly writing data into an optic disk that is performed by an optic disk drive incorporating a control unit to which an external data storage device, such as compact flash memory device, is connected. The method comprises steps of (1) initiating a writing operation, (2) setting the optic disk drive to busy condition, (3) checking if an optic disk is properly loaded and if the external memory device is correctly connected, (4) checking if the optic disk is a UDF disk; (5) issuing a warning, if it is not, (6) creating a folder in the optic disk, (7) retrieving data from the external data storage device and writing the data into the folder of the optic disk, and (8) ending the writing operation. No computer-based interface is required between the optic disk drive and the external data storage device in performing the data writing operation."}
-{"text": "This invention relates to a synchronization signal processing system for use in a mobile communication network which comprises a plurality of mobile service switching centers and a plurality of base transceiver stations and is operable in a time division fashion.\nThe mobile communication network has an overall service area which is divided into cells or radio zones assigned with the base transceiver stations, respectively, and in which a plurality of mobile stations are present, namely, either moving or staying standstill, at a time. Each mobile station may be either a portable telephone device carried by a user or a subscriber's terminal installed in an automobile or in a like mobile vehicle and is movable from a first zone of the cells to a second zone of the cells.\nIt is possible to understand that each mobile service switching center is connected to a plurality of fixed subscriber substations either directly or through at least one exchange office. Some of the mobile service switching centers are connected to the base transceiver stations. More particularly, each of such mobile service switching centers is connected to a certain number of base transceiver stations.\nThe mobile service switching centers are connected to one another by wired communication lines. The mobile service switching centers and the base transceiver stations may be connected through wired communication lines. Among the overall service area, some of the cells are often referred to collectively as a radio communication area when assigned to the base transceiver stations which are served by one of the mobile service switching centers.\nEach base transceiver station is for transmitting and receiving radio message signals to and from at least one of the mobile stations that is currently present in the cell assigned with the base transceiver station under consideration. For use in time division multiple access (TDMA), the radio message signals are carried by a radio carrier signal of a radio frequency in a plurality of time slots. A predetermined number of such time slots are successively arranged in a frame in the manner known in the art.\nWhen a particular station of the mobile stations moves between the first and the second zones assigned with first and second stations of the base transceiver stations, the first and the second stations use different radio frequencies and different time slots in transmitting and receiving the radio message signals to and from the particular station. The first and the second stations may be connected either to one or to two of the base transceiver stations. In either event, the particular station is inevitably subjected to a handover processing between the first and the second stations. It is therefore desirable to preliminarily synchronize the frames and the time slots in the base transceiver stations in order to reduce a time necessary for such a handover processing as a handover processing time.\nIn the manner which will later be described, a conventional synchronization signal processing system comprises an individual synchronization signal generating circuit in each mobile service switching center. When connected to such a mobile service switching center, the base transceiver station can generate synchronized frames and synchronized time slots for the mobile stations which are currently present in the radio communication area served by the base transceiver station under consideration.\nA little more in detail, the synchronization signal generating circuit comprises first and second time division switches, each comprising controllable connection paths and producing a switch trouble signal when a trouble occurs therein. A controller device is cross connected to the first and the second time division switches and is supplied with the switch trouble signal to control the connection paths of one of the first and the second time division switches that is not producing the switch trouble signal and serves as an active switch with the other of the first and the second time division switches used as a standby switch. A synchronization signal generator is connected to the active switch to supply a synchronization signal to the connection paths of the active switch. Output trunk circuits are connected to the connection paths of the first and the second time division switches to supply the synchronization signal to at least one of the output trunk circuit from the connection paths controlled by the controller device to the base transceiver stations served by mobile service switching center in question.\nIt is liable that the synchronization signal generator is involved into a trouble. First and second synchronization signal generators are therefore cross connected to the first and the second time division switches. Alternatively, it is possible to understand that the first and the second synchronization signal generators are connected to the active switch. In either event, each synchronization signal generator produces a generator trouble signal when a trouble occurs therein. Supplied with the generator trouble signal, the controller device controls the connection paths of the active switch to supply the output trunk circuits with the synchronization signal generated by one of the first and the second synchronization signal generators that is not producing the generator trouble signal.\nAs a consequence, the conventional synchronization signal processing system can deal with troubles that may occur in the time division switches and/or in the synchronization signal generators. It is, however, impossible to keep the phase of the synchronization signal when the first and the second synchronization signal generators are switched from one to the other."}
-{"text": "Arc discharge in aqueous electrolytes (for example, welding under seawater), is widely used in engineering and construction, and is at present the only known form of stationary plasma discharge in liquid media. In recent years, such discharge was also used in different physicochemical studies and in the synthesis of various materials. The specific feature of arc discharge in liquid media is the localization of a plasma region near the electrode ends and a \u201cfalling\u201d form of volt-ampere characteristic as illustrated in FIG. 1.\nIn a gaseous phase, different kinds of discharges can be implemented, the external manifestation and electrical parameters of which are connected with a wide range of technical characteristics for devices used in their implementation and a variety of elementary processes determining the conditions of current passage through gas. The essential feature of electric discharge development in the gaseous phase is a profound effect of the properties of the gas medium on the current passages through the gas.\nUnder usual conditions, the concentration of charge carriers (electrons and ions) in the gas is very low: a gas is a very good dielectric. For a gas to have a high electrical conductivity (as a result of ionization) it is necessary for a high quantity of charge carriers to be present, requiring in turn a great quantity of energy. Gases have a steady electric conductivity when there is equilibrium between the origination and disappearance of charges. Thus, to create a means by which high electrical conductivity in a gas can be achieved through substantially lower energy requirements than has been taught in the prior art is highly desirable.\nIf the rate of movement of electrical charges is proportional to the field strength, the conductivity of gas approximately obeys Ohm's law (FIG. 2, section a). With increasing field strength, the decrease of electrical charges begins to have an influence (FIG. 2, section b) because of the migration of the charges to the electrodes. Further increases of the electrical field strength result in a steep increase of current due to the start of collision ionization (FIG. 2, section c). In spite of the avalanche-like character of current increases, the existence of external ionizer(s) is needed to sustain the electrical discharge, and the discharge remains being as not self-sustained (region 1). Eventually, a point is reached where for each electron leaving the cathode, one or more electrons arrive at the anode, in a phenomenon known as breakdown discharge (glow discharge or plasma discharge). This causes a self-sustained electrical current from the cathode to the anode. However, the current state-of-the-art process requires a large amount of energy to reach this self-sustaining threshold. Since high energy requirements directly and indirectly decrease the overall economy of the model, the requirement of high energy is undesirable. Therefore, it is highly desirable to have a new process having low energy demand in which the transition from non self-sustained discharge to self-sustained discharge (glow discharge) would occur with a low-energy input.\nAgain referring to FIG. 1, which illustrates the prior art, the voltage-current characteristic curve for glow discharge preferably comprises three sections, referred to for the sake of clarity as subnormal section or subnormal mode (FIG. 2, section d), normal section or normal mode (FIG. 2, section e) and abnormal section or abnormal mode (FIG. 2, section f).\nFurther increase of current density on the cathode causes the appearance of electric arc, as well as a drastic change of the main characteristics of the discharge (FIG. 2, section g).\nIt should be noted that the appearance or threshold of discharges in the gas phase depends considerably on the pressure of the gas. Thus, in the case of a uniform field of breakdown voltage (self-maintained discharge initiation voltage) the threshold is determined by the product of pressure by the distance between the electrodes, according to Paschen's Law. Pachen determined that breakdown voltage is determined by the following equation:\n  V  =            a      \u2061              (        pd        )                            ln        \u2061                  (          pd          )                    +      b      where V is the breakdown voltage in Volts, p is the pressure in atmospheres, d is the gap distance in meters, and a and b are constants that depend upon the particular gas between the electrodes. Thus, in contrast to liquids, which are relatively incompressible, different forms of electric discharge can be implemented in gases by varying the pressure of the gas between the electrodes.\nMoreover, when ultrasonic cavitation, a sort of \u201ccold boiling\u201d resulting from the creation and collapse of zillions of microscopic bubbles in the liquid caused by ultrasonic waves, is implemented within a liquid, its phase composition and physical properties abruptly change, which can lead to some specific features for the formation of electric discharges within the liquid. In the region of intense cavitation, a gaseous component is formed which represents a significant fraction of the liquid. Therefore it can be assumed that the conditions for electric breakdown into the cavitation region should become easier, and the initiation of different forms of discharge could start through use of this invention. By varying the parameters of an ultrasonic field, it is possible to influence the processes of plasma glow within a cavitating liquid.\nThe prior art has several examples of attempts to resolve this problem.\nHowever, few patent applications or patents work in the abnormal mode. In abnormal mode, also known as abnormal glow, effectively all of the gas molecules must be ionized to provide charge carriers for the current. Typically, the gas molecules are ionized multiple times meaning that more than one electron has been freed for most of the gas molecules. This creates a relatively uniformly distributed plasma across the electrodes. A higher density (or pressure) of gas molecules, on the other hand, would lead to a normal mode, or normal glow discharge. In this region, fewer than all of the molecules are ionized. This creates a situation where plasma forms in a relatively small region between the electrodes. A plasma discharge of this type can lead to concentrated energy in a relatively small area and possibly lead to electrode damage. Therefore, it is preferable to work in the abnormal mode.\nThose patent applications or patents that do work in the abnormal mode, like U.S. Pat. No. 5,068,002, to Monroe, do not use an electrode as the radiator, in the same way that the instant application uses it, whereby the current application discloses a very low energy consumption jointly with a very low voltage to initiate and maintains a volumetric discharge which generates operational advantages in term of achieving the goals of this application. Monroe describes an ultrasonic glow discharge surface cleaning apparatus for abrading contaminants from the surface of a work piece using plasma glow discharge.\nFor example, in US Patent Application 2004/0265137 A1 to Bar-Gadda, a method is proposed for hydrogen production from water or steam by means of plasma discharge excited in the UHF, radio- or low-frequency range, as well as with arc discharge. This application describes the injection of water molecules into plasma discharge.\nU.S. Pat. No. 7,070,634 B1 A1 to Wang describes a plasma apparatus for converting a gaseous mixture of water vapor and hydrocarbons into hydrogen.\nUS Patent Application 2006/0060464 to Chang teaches a fluid phase contained in a reactor, within which electrodes (anode and cathode) are placed. A flow of gas bubbles is introduced or generated in the medium in the region adjacent to the cathode. The potential difference necessary for the initiation of glow discharge and for the ionization of gas molecules in the bubbles is applied between the cathode and the anode.\nU.S. Pat. No. 7,067,204 to Nomura et al., describes an apparatus comprising an ultrasonic generator for creation of bubbles within a liquid, and a generator providing the excitation of electromagnetic waves in the liquid phase, for the implementation of the plasma discharge.\nJapanese Application JP2006273707 to Shibata et al. relates to the publication, \u201cSynthesis of amorphous carbon nanoparticles and carbon-encapsulated metal nanoparticles in liquid benzene by an electric plasma discharge in ultrasonic cavitation field,\u201d Ultrasonic Sonochemistry 13 (2006) 6-12, Institute of Multidisciplinary Research for Advanced Material (IMRAM), Tohoku University. This application illustrates a method and a device for producing a nanocarbon material that does not require an expensive production facility such as the ones normally required for dry treatment. It can easily produce the nanocarbon material because the application of high voltage is not needed and neither worsens nor deteriorates the working environment in a production premise, and at the same time considers safety factors. This method can remarkably reduce production costs by improving production efficiency because of its continuous production and recovery, and providing an alternative for mass productivity. The method comprises a process (A) for arranging electrodes, one cathode and one anode, connected to the power source; an ultrasonic horn connected to an ultrasonic generator within an organic solvent that fills a container; and a process (B) for generating an ultrasonic cavitation field by ultrasonic waves into the organic solvent, around the head of the ultrasonic horn; and effecting the thermal decomposition of the molecules in the organic solvent by applying a voltage to the electrodes so as to generate plasma discharge within the ultrasonic cavitation field adequate for the production of the nanocarbon material.\nU.S. Pat. No. 6,835,523 to Yamazaki et al. describes a \u201cMethod for fabricating with ultrasonic vibration a carbon coating,\u201d which is a process for fabricating a carbon coating in a medium disposed on one side of an electrode connected to a high-frequency power supply. Ultrasonic vibrations are then supplied to the object.\nNone of the prior art, however, either individually or in combination, provides a method by which initiating and maintaining an abnormal glow volumetric sonoplasma discharge can be performed using a substantially lower amount of electrical power.\nThus there has existed a long-felt need for a method by which the sonoplasma discharge can be initiated and maintained with substantially less electrical power than is currently needed to accomplish the same result using the prior art. This is accomplished with this invention.\nThe current invention provides just such a solution by having a method and apparatus for initiating and maintaining an abnormal glow volumetric sonoplasma discharge (VSPD). With certain parameters of the electrical discharge and of the intensity of elastic vibrations, it is possible to initiate VSPD within a cavitating liquid medium. The mechanism for the initiation of VSPD is related to the breakdown of gas-phase microchannels formed by the growth cavitation bubbles. The method uses elastic vibrations (EV) in the frequency range 1,000-100,000 Hz with enough intensity for the development of cavitation phenomena; these vibrations are introduced into the liquid-phase working medium, and a source of direct, alternating (hertz and kilohertz range), high frequency (HF) (megahertz range) and ultrahigh frequency (UHF) (gigahertz range) electric field in liquid (DPS) provides the initiation and stable glow of VSPD. Resulting VSPD is characterized by volumetric glow in the frequency range of visible light and ultraviolet radiation in the entire cavitation-electric field, and is characterized by a rising volt-ampere characteristic curve.\nWhen a high-intensity ultrasonic field exceeding a cavitation threshold is induced within liquids, a new form of electric discharge is obtained, characterized by a volumetric glow electrical discharge throughout the space between the electrodes, having a rising volt-ampere characteristic curve that is inherent to abnormal glow discharge in gas. Such discharge within the liquid has the surface characteristic of micro bubbles, and can be used for the design of novel sonoplasma-chemical processes because of the extensive interface plasma. The heterogeneous liquid/gas-vapor system leads to a rise in diffusion rates of chemically active particles in the system and a more economical method to achieve the desired result(s)."}
-{"text": "Over the last two decades, mass spectrometry has made tremendous strides in analyzing protein samples derived from a variety of different sample types. Coupled with electrospray ionization and various separation techniques, thousands of proteins may be identified and quantitated in a single sample. The most common approach used in the laboratory today involves some form of protein extraction followed by proteolytic digestion of protein sample of interest. The use of proteolytic enzymes like trypsin produces peptides that can easily be analyzed by a variety of different instrument configurations. This approach termed \u201cbottom-up\u201d proteomics, can be used to study the state of living cells as a function of their environment. One of the major advantages of the \u201cbottom-up\u201d approach is that the peptides produced have very similar physiochemical properties which makes for a straight forward separation of thousands of peptides in complex samples. Any separation approach coupled with tandem mass spectrometry can then be used to produce amino acid sequence information that is utilized to identify the proteins in a given sample. Although this technique is routine in many laboratories, there are limitations as to the amount of information that can be obtained when reducing intact proteins to their constituent peptides.\nIn contrast to \u201cbottom-up\u201d proteomics, \u201ctop-down\u201d proteomics refers to methods of analysis in which protein samples are introduced intact into a mass spectrometer, without enzymatic, chemical or other means of digestion. Top-down analysis enables the study of the intact protein, allowing identification, primary structure determination and localization of post-translational modifications (PTMs) directly at the protein level. Top-down proteomic analysis typically consists of introducing an intact protein into the ionization source of a mass spectrometer, fragmenting the protein ions and measuring the mass-to-charge ratios and abundances of the various fragments so-generated. The resulting fragmentation is many times more complex than a peptide fragmentation, which may, in the absence of the methods taught herein, necessitate the use of a mass spectrometer with very high mass accuracy and resolution capability in order to interpret the fragmentation pattern with acceptable certainty. The interpretation generally includes comparing the observed fragmentation pattern to either a protein sequence database that includes compiled experimental fragmentation results generated from known samples or, alternatively, to theoretically predicted fragmentation patterns. For example, Liu et al. (\u201cTop-Down Protein Identification/Characterization of a Priori Unknown Proteins via Ion Trap Collision-Induced Dissociation and Ion/Ion Reactions in a Quadrupole/Time-of-Flight Tandem Mass Spectrometer\u201d, Anal. Chem. 2009, 81, 1433-1441) have described top-down protein identification and characterization of both modified and unmodified unknown proteins with masses up to \u224828 kDa\nAn advantage of a top-down analysis over a bottom-up analysis is that a protein may be identified directly, rather than inferred as is the case with peptides in a bottom-up analysis. Another advantage is that alternative forms of a protein, e.g. post-translational modifications and splice variants, may be identified. However, top-down analysis has a disadvantage when compared to a bottom-up analysis in that many proteins can be difficult to isolate and purify. Thus, each protein in an incompletely separated mixture can yield, upon mass spectrometric analysis, multiple ion species, each species corresponding to a different respective degree of protonation and a different respective charge state, and each such ion species can give rise to multiple isotopic variants.\nThe process of analyzing intact proteins in cell lysates by mass spectrometry (MS) is associated with a number of difficulties. Firstly, electrospray ionization (ESI) of protein mixtures from cell lysates can generate extremely complex mass spectra due to the presence of multiple proteins, each comprising its own charge state envelope, where each charge state envelope is the collection of mass spectral lines corresponding to plural charge states, and where each charge state correlates directly with the number of positively charged protons that are adducted to an otherwise charge-free molecule. Consequently, multiple charge state envelopes may be overlapping within any given mass-to-charge (m/z) range. In this example, multiple proteins overlap at the same m/z value that have different molecular weights and charges. Commonly used techniques in MS are often insufficient for simplifying these spectra because of the inherent peak overlapping as well as the inherent wide range of magnitudes of MS lines of ionized constituents, where such constituents may range from uninteresting small molecules to interfering biomolecules to the proteins of interest, themselves. Isolation of a specified charge state of a protein within such complex spectra does not typically alleviate the burden of multiple protein peaks overlapping, since the isolation of ions of a particular protein charge state will generally result in co-isolation of one or more additional ions. This co-isolation makes it a challenge not only to dissociate the protein in an attempt to identify it based on the fragments produced, but also to accurately determine the intact mass and sequence coverage of that protein.\nSo-called \u201cfront-end\u201d separation techniques, such as liquid chromatography (LC) or ion mobility spectrometry (IMS), performed prior to introduction of samples into a mass spectrometer, may be implemented to reduce the overall complexity and provide an additional major benefit, which is the reduction of ionization competition at an ionization source. Unlike mixtures of proteolytic peptides typically analyzed in bottom-up experiments, intact proteins mixtures contain a wide range of molecular weights, isoelectric points, hydrophobicities, and other physiochemical properties that make it challenging to analyze these mixtures via any single separation technique in a comprehensive manner. Both of the above separation methods are associated with their own benefits and pitfalls. Liquid chromatography tends to require significant amounts of time per sample to separate individual proteins, although it is still common to have two or more proteins co-elute. Enhanced separation can reach the point of becoming more of \u201can art\u201d than a standardized method, and the enhanced separation may be dependent on the user skill in the state-of-the-art. The latter technique, IMS, can rapidly separate certain proteins and/or charge states from others but IMS spectra are at least partially correlative with (i.e., not \u201corthogonal to\u201d) mass spectra. The IMS method also suffers from ionization competition, requires extensive optimization and typically involves dynamic conditions to observe a full mass spectrum containing all charge states.\nProton transfer reactions, a type of ion-ion reaction that has been used extensively in biological applications for rapid separations of complex mixtures, addresses many of these aforementioned concerns. Experimentally, proton transfer is accomplished by causing multiply-positively-charged protein ions from a sample to react with introduced singly-charged reagent anions so as to reduce the charge of the multiply-charged protein ions. These reactions proceed with pseudo-first order reaction kinetics when the anions are present in large excess over the protein ion population. The rate of reaction is directly proportional to the square of charge of the protein ion (or other multiply-charged cation) multiplied by the charge on the anion. The same relationship holds for reactions of the opposite polarity as well. This produces a series of pseudo-first order consecutive reaction curves as defined by the starting multiply-charged protein ion population. Although the reactions are highly exothermic (in excess of 100 kcal/mol), proton transfer is an even-electron process performed in the presence of 1 mtorr of background gas (i.e. helium) and thus does not fragment the starting multiply-charged protein ion population. The collision gas serves to remove the excess energy on the microsecond time scale (108 collisions per second), thus preventing fragmentation of the resulting product ion population.\nProton transfer reactions (PTR) have been used successfully to identify individual proteins in mixtures of proteins. This mixture simplification process has been employed to determine charge state and molecular weights of high mass proteins. PTR has also been utilized for simplifying product ion spectra derived from the collisional-activation of multiply-charged precursor protein ions. Although PTR reduces the overall signal derived from multiply-charged protein ions, this is more than offset by the significant gain in signal-to-noise ratio of the resulting PTR product ions. The PTR process is 100% efficient leading to only single series of reaction products, and no side reaction products that require special interpretation and data analysis.\nVarious aspects of the application of PTR to the analysis of peptides, polypeptides and proteins have been described in the following documents: U.S. Pat. No. 7,749,769 B2 in the names of inventors Hunt et al., U.S. Patent Pre-Grant Publication No. 2012/0156707 A1 in the names of inventors Hartmer et al., U.S. Pre-Grant Publication No. 2012/0205531 A1 in the name of inventor Zabrouskov; McLuckey et al., Anal. Chem. 1998, 70:1198-1202; Stephenson et al., J. Am. Soc. Mass Spectrom. 1998, 8:637-644; Stephenson et al., J. Am. Chem. Soc. 1996, 118:7390-7397; McLuckey et al., Anal. Chem. 1995, 67:2493-2497; Stephenson et al., Anal. Chem. 1996, 68:4026-4032; Stephenson et al., J. Am. Soc. Mass Spectrom. 1998, 9:585-596; Stephenson et al., J. Mass Spectrom. 1998, 33:664-672; Stephenson et al., Anal. Chem., 1998, 70:3533-3544 and Scalf et al., Anal. Chem. 2000, 72:52-60. Various aspects of general ion/ion chemistry have been described in McLuckey and Stephenson, Mass Spec Reviews 1998, 17:369-407 and U.S. Pat. No. 7,550,718 B2 in the names of inventors McLuckey et al. Apparatus for performing PTR and for reducing ion charge states in mass spectrometers have been described in U.S. Pre-Grant Publication No. 2011/0114835 A1 in the names of inventors Chen et al., U.S. Pre-Grant Publication No. 2011/0189788 A1 in the names of inventors Brown et al., U.S. Pat. No. 8,283,626 B2 in the names of inventors Brown et al. and U.S. Pat. No. 7,518,108 B2 in the names of inventors Frey et al. Adaptation of PTR charge reduction techniques to detection and identification of organisms has been described by McLuckey et al. (\u201cElectrospray/Ion Trap Mass Spectrometry for the Detection and Identification of Organisms\u201d, Proc. First Joint Services Workshop on Biological Mass Spectrometry, Baltimore, Md., 28-30 Jul. 1997, 127-132).\nThe product ions produced by the PTR process can be accumulated into one or into several charge states by the use of a technique known as \u201cion parking\u201d. Ion parking uses supplementary AC voltages to consolidate the PTR product ions formed from the original variously protonated ions of any given protein molecule into a particular charge state or states at particular m/z values during the reaction period. This technique can be used to concentrate the product ion signal into a single or limited number of charge states (and, consequently, into a single or a few respective mass-to-charge [m/z] values) for higher sensitivity detection or further manipulation using collisional-activation, ETD, or other ion manipulation techniques. Various aspects of ion parking have been described in U.S. Pat. No. 8,440,962 B2 in the name of inventor Le Blanc and in the following documents: McLuckey et al., Anal. Chem. 2002, 74:336-346; Reid et al., J. Am. Chem. Soc. 2002, 124:7353-7362; He et al., Anal. Chem. 2002, 74:4653-4661; Xia et al., J. Am. Soc. Mass. Spectrom. 2005, 16:71-81; Chrisman et al., Anal. Chem. 2005, 77:3411-3414 and Chrisman et al., Anal. Chem. 2006, 78:310-316.\nAnother difficulty associated with the mass spectrometric analysis of proteins in cell lysates by (MS) is that the fragmentation behavior for each charge state of a protein is generally unknown prior to the dissociation event. In particular, ions comprising some charge states can dissociate well while ions comprising other charge states may dissociate poorly. Isolation and dissociation of ions of a particular charge state therefore does not guarantee efficient dissociation or dissociation into a set of diagnostic fragments.\nA third challenge associated with intact protein analysis is the wide distribution of charge states produced for high molecular weight proteins typically in excess of 50 kDa. Here the starting signal can be divided into over 30 plus charge states, making tandem mass spectrometry of any given charge state produce a spectrum with low signal-to-noise ratio. The ability to produce ample sequence coverage for protein identification can therefore be difficult with a single tandem mass spectrum.\nA variety of ion activation (fragmentation) techniques can be used to produce structural information on intact proteins. The most commonly used approach termed collision-induced dissociation (CID) involves collisions of an isolated population of multiply-charged precursor ions with a neutral background gas. Most commonly, the multiply-charged precursor ions are accelerated using the fundamental frequency of motion of the defined ion population in order to collide with the neutral background gas so as to produce unimolecular dissociation events. This process leads to fragmentation along the amide backbone of the protein thus yielding amino acid sequence information. More extensive fragmentation of proteins can be obtained with higher collision energy processes termed HCD or high energy collision induced dissociation. Many times this involves multiple fragmentation events inside the collision cell thus producing more extensive sequence coverage. Another approach used to produce protein sequence coverage via ion activation is that of photodissociation (PD), where photons of a defined wavelength are used to excite the ion of interest. Two common types employed include ultra-violet (UV-PD) and infrared multiphoton dissociation (IRMPD). The latter is a high energy process where the rate of energy deposition in the ion far exceeds that of the dissociation process. Here fragmentation can be produced along any point in the protein backbone, or may yield amino acid side chain fragmentation as well. For IRMPD, this is a much lower energy process that is characterized by the presence of cleavages at amide bonds and losses of ammonia and water from the intact protein and fragment ions generated during irradiation. The time frame of the IRMPD experiment can be expanded to produce more extensive fragmentation as well. Ion-ion reactions using electron transfer reagent ions can also be employed as a fragmentation approach for intact proteins. Here an electron transfer event from the multiple-charged protein to the singly-charged anion produces backbone fragmentation of the protein with any posttranslational modifications still intact.\nTaken together, these ion activation approaches for tandem mass spectrometry produce many different complementary forms of fragmentation that can provide protein sequence information. Ideally, these approaches can be applied in a broad band fashion in order to increase sequence coverage of proteins and provide additional information on modifications, splice variants, and expression of single amino acid mutations. The application of these approaches in a broadband format (i.e. covering multiple charge states of the same intact protein) would provide a more comprehensive view of protein characterization and identification."}
-{"text": "1. Field of the Invention\nThe present invention relates to a data recording apparatus for recording data on a record (recording, recordable or recorded) medium, a method therefor, a data reproducing apparatus for reproducing data recorded on a record medium, a method therefor and a record medium on which data has been recorded. More particularly, the present invention relates to an information providing/collecting apparatus for providing and collecting so-called multimedia information, such as video information and music information, or program information and a method therefor.\n2. Related Background Art\nAs a data record medium on which information signals, such as Audio data, video data and various data items, are recorded, means for optically recording information signals, specifically, a so-called compact disk (CD) for use in the music field and a CD-ROM which meets the CD standard and which is used for data have been used all over the world in recent years.\nHitherto, information providing service has been realized as a so-called data base system and a personal computer communication system in each of which a user terminal (a terminal of an information collecting side) and an information provider are connected to each other through, for example, the telephone line to enable information required by the user to be taken out. Another information providing service has been realized with which a large-capacity medium, such as a so-called CD-ROM having encoded information recorded thereon is distributed and key information for decoding encoded information is transmitted to the user by, for example, communication so that encoded information recorded on the CD-ROM is decoded and decoded information is copied on a hard disk or the like so as to be used.\nMoreover, a technique has been disclosed in Japanese Patent Publication No. 2-60007 in which a password formed by encoding a file key by using a code key is supplied to a computer; and a program written on the record medium is decoded by a coding mechanism to prevent copying and sharing of the software program.\nHitherto, all of information items recorded on the foregoing CD or the CD-ROM are read by a reproducing apparatus and copied onto, for example, a hard disk. Then, data copied onto the hard disk is supplied to an encoder system for the CD or the CD-ROM to newly make a CD or a CD-ROM so that a pirate edition is easily manufactured. As described above, the security function, such as the copy protection, has been unsatisfactory.\nThe foregoing problem is also critical for a so-called digital video disk (DVD), which is expected to be a data record medium for a next generation.\nOn the other hand, in the conventional information providing service, a method has been employed in which key information for decoding is transmitted to a user in such a manner that key information is transmitted by means of voice through a telephone line. Thus, key information has not been encoded particularly. However, the foregoing method has a risk in view of keeping security.\nIn the case where communication is employed to transmit key information, one-to-one connection is usually established. Therefore, there is substantially no risk of key information being stolen. However, in the case where key information is transmitted through a network, there arises a problem in protecting key information.\nTherefore, in an information providing system, in which mediums, on each of which encoded information has been recorded in a large quantity, are distributed by the information provider; and only in a case where a user requires information to obtain from the medium, key information for decoding the code is supplied and accounting is performed, the problem in view of security when key information is transmitted results in a risk to arise in that key information can be obtained by a person except the subject user. In the foregoing case, the information providing system cannot be held. If whether or not the user is a formal user cannot be specified, there is a risk that account is put down to'another person. Also in the foregoing case, the information providing system cannot be held.\nThus, security improvement in transmitting key information from an information provider to a user and reliable specification of a user are important requirements."}
-{"text": "1. Field of the Invention\nThis invention is directed to a unique method and device for delivering controlled heat to perform ablation to treat benign prosthetic hypertrophy or hyperplasia (BPH). The method and the apparatus deliver this controlled heat into tissue penetrated by devices such as those disclosed in the copending above-referenced applications.\n2. Discussion of Background\nTreatment of cellular tissues usually requires direct contact of target tissue with a medical instrument, usually by surgical procedures exposing both the target and intervening tissue to substantial trauma. Often, precise placement of a treatment probe is difficult because of the location of a target tissue in the body or the proximity of the target tissue to easily damaged, critical body organs, nerves, or other components.\nBenign prostatic hypertrophy or hyperplasia (BPH), for example, is one of the most common medical problems experienced by men over 50 years old. Urinary tract obstruction due to prostatic hyperplasia has been recognized since the earliest days of medicine. Hyperplastic enlargement of the prostate gland often leads to compression of the urethra, resulting in obstruction of the urinary tract and the subsequent development of symptoms including frequent urination, decrease in urinary flow, nocturia, pain, discomfort, and dribbling. The association of BPH with aging has been shown to exceed 50% in men over 50 years of age and increases in incidence to over 75% in men over 80 years of age. Symptoms of urinary obstruction occur most frequently between the ages of 65 and 70 when approximately 65% of men in this age group have prostatic enlargement.\nCurrently there is no proven effective nonsurgical method of treatment of BPH. In addition, the surgical procedures available are not totally satisfactory. Currently patients suffering from the obstructive symptoms of this disease are provided with few options: continue to cope with the symptoms (i.e., conservative management), submit to drug therapy at early stages, or submit to surgical intervention. More than 30,000 patients per year undergo surgery for removal of prostatic tissue in the United States. These represent less than five percent of men exhibiting clinical significant symptoms.\nThose suffering from BPH are often elderly men, many with additional health problems which increase the risk of surgical procedures. Surgical procedures for the removal of prostatic tissue are associated with a number of hazards including anesthesia associated morbidity, hemorrhage, coagulopathies, pulmonary emboli and electrolyte imbalances. These procedures performed currently can also lead to cardiac complications, bladder perforation, incontinence, infection, urethral or bladder neck stricture, retention of prostatic chips, retrograde ejaculation, and infertility. Due to the extensive invasive nature of the current treatment options for obstructive uropathy, the majority of patients delay definitive treatment of their condition. This circumstance can lead to serious damage to structures secondary to the obstructive lesion in the prostate (bladder hypertrophy, hydronephrosis, dilation of the kidney pelves, etc.) which is not without significant consequences. In addition, a significant number of patients with symptoms sufficiently severe to warrant surgical intervention are poor operative risks and are poor candidates for prostatectomy. In addition, younger men suffering from BPH who do not desire to risk complications such as infertility are often forced to avoid surgical intervention. Thus the need, importance and value of improved surgical and non-surgical methods for treating BPH is unquestionable.\nHigh-frequency currents are used in electrocautery procedures for cutting human tissue especially when a bloodless incision is desired or when the operating site is not accessible with a normal scalpel but presents an access for a thin instrument through natural body openings such as the esophagus, intestines or urethra. Examples include the removal of prostatic adenomas, bladder tumors or intestinal polyps. In such cases, the high-frequency current is fed by a surgical probe into the tissue to be cut. The resulting dissipated heat causes boiling and vaporization of the cell fluid at this point, whereupon the cell walls rupture and the tissue is separated.\nDestruction of cellular tissues in situ has been used in the treatment of many diseases and medical conditions alone or as an adjunct to surgical removal procedures. It is often less traumatic than surgical procedures and may be the only alternative where other procedures are unsafe. Ablative treatment devices have the advantage of using a destructive energy which is rapidly dissipated and reduced to a non-destructive level by conduction and convection forces of circulating fluids and other natural body processes.\nMicrowave, radiofrequency, acoustical (ultrasound) and high energy (laser) devices, and tissue destructive substances have been used to destroy malignant, benign and other types of cells and tissues from a wide variety of anatomic sites and organs. Tissues treated include isolated carcinoma masses and, more specifically, organs such as the prostate, glandular and stromal nodules characteristic of benign prostate hyperplasia. These devices typically include a catheter or cannula which is used to carry a radiofrequency electrode or microwave antenna through a duct to the zone of treatment and apply energy diffusely through the duct wall into the surrounding tissue in all directions. Severe trauma is often sustained by the duct wall during this cellular destruction process, and some devices combine cooling systems with microwave antennas to reduce trauma to the ductal wall. For treating the prostate with these devices, for example, heat energy is delivered through the walls of the urethra into the surrounding prostate cells in an effort to kill the tissue constricting the urethra. Light energy, typically from a laser, is delivered to prostate tissue target sites by \"burning through\" the wall of the urethra. Healthy cells of the duct wall and healthy tissue between the nodules and duct wall are also indiscriminately destroyed in the process and can cause unnecessary loss of some prostate function. Furthermore, the added cooling function of some microwave devices complicates the apparatus and requires that the device be sufficiently large to accommodate this cooling system.\nApplication of liquids to specific tissues for medical purposes is limited by the ability to obtain delivery without traumatizing intervening tissue and to effect a delivery limited to the specific target tissue. Localized chemotherapy, drug infusions, collagen injections, or injections of agents which are then activated by light, heat or chemicals would be greatly facilitated by a device which could conveniently and precisely place a fluid supply catheter opening at the specific target tissue."}
-{"text": "Unless otherwise indicated herein, the materials described in this section are not prior art to the claims in this disclosure and are not admitted to be prior art by inclusion in this section.\nA news production system (NPS) may facilitate the production of a news program in the form of a media stream. In one example, an NPS may include multiple media sources and a production switcher, where outputs of the media sources are connected to inputs of the production switcher. This may allow the production switcher to switch between and/or combine multiple media streams output by the media sources, thereby outputting the news program in the form of another media stream.\nThere are various types of media, including for example, audio, video, or a combination thereof. As such, in one example, an NPS may output a news program in the form of an audio stream. In this instance, the NPS may transmit the audio stream to a radio-broadcasting system for broadcast. As another example, a media stream may take the form of a video stream or a combined audio and video stream. In such instances, the NPS may transmit the video stream or the combined audio and video stream to a television-broadcasting system for broadcast.\nA media source may take a variety of forms. For example, a media source may take the form of a media server. A media server is a device configured for retrieving a media file, converting the retrieved media file into a media stream, and outputting the converted media stream.\nAs another example, a media source may take the form of a media effect engine. A media effect engine is a device configured for retrieving a media effect (sometimes referred to as a \u201cpage\u201d), and running the media effect thereby outputting a corresponding media stream. A media effect may be stored as a file that includes instructions and other data (e.g., media) related to the media effect. By running the media effect, the media effect engine may generate and output a media stream based on those instructions. Media effects are commonly used as a means to generate animations, graphics, or other visual effects in the form of a media stream that can be overlaid on another media stream. For instance, in the context of a news program, a \u201clower third\u201d media effect may be used to overlay a graphic over a lower third portion of a media stream.\nAs such, in one example NPS, a media server may output a first media stream while a media effect engine outputs a second media stream, and a production switcher may combine the two media streams (e.g., by overlaying the second media stream over the first media stream) to output the news program in the form of a third media stream. The production switcher may then transmit the third media stream to a broadcasting system (e.g., a television-broadcasting system) for broadcast.\nA media effect engine may be controlled in a variety of ways such that it may perform the steps of retrieving a media effect and running the media effect. For instance, a user may control a media effect engine by providing it a suitable instruction via a user-interface. However, for a variety of reasons, this manner of controlling the media effect engine may be undesirable. Among other things, this process may be time-consuming for the user. In addition, it may be difficult for the user to ensure that the media effect engine performs such steps at appropriate times during production of the news program.\nAs another example, a controller device may control a media effect engine by providing it a suitable instruction in accordance with one or more application programming interfaces (API) that may be made publically available by the provider of the media effect engine or another entity. However, again for a variety of reasons, this manner of controlling the media effect engine may be undesirable. Among other things, it may be time-consuming for a user to configure the controller device to provide such a suitable instruction.\nThis approach may be particularly time-consuming given that different instructions may need to adhere to different APIs. For instance, an instruction requesting the running of one type of media effect may need to adhere to a different API than another instruction requesting the running of another type of media effect. In addition, it may be difficult for the user to configure the controller device such that it causes the media effect engine to retrieve and run a particular media effect at an appropriate time during production of the news program."}
-{"text": "1. Field of the Invention\nThis invention relates to a dry etching method adapted for fine processing of manufacturing a semiconductor device, and particularly to a method for preventing regression of a resist mask formed on an SiON based antireflection film so as to improve anisotropy.\n2. Description of Related Art\nIn order to realize large scale integration of semiconductor devices, the minimum processing size of the circuit pattern formation has been rapidly diminished. For instance, the minimum processing size of the 16M DRAM of approximately 0.5 .mu.m (half micron), the minimum processing size of the 64M DRAM of 0.35 .mu.m (sub-half micron), and the minimum processing size of the 256M DRAM of 0.25 .mu.m (quarter micron) are required.\nThis increasingly fine processing depends largely upon a technique of photolithography to form a mask pattern. Visible to near ultraviolet rays, such as g rays having a wavelength of 436 nm or i rays having a wavelength of 365 nm, of a high pressure mercury lamp are used for the current 0.5-.mu.m class processing, and far ultraviolet rays, such as KrF excimer laser lights having a wavelength of 248 nm, are used for 0.35 to 0.25-.mu.m class processing. In the photolithography technique for forming a fine mask with a ray width of not greater than 0.4 .mu.m, an antireflection film to weaken a reflected light from an underlying material layer is substantially required for preventing reduction in contrast and resolution due to halation and standing wave effect.\nAs the component material of the antireflection film, amorphous silicon, TiN and TiON are conventionally used. However, since it has been shown that SiON (silicon oxide nitride) exhibits satisfactory optical properties in the far ultraviolet region, application of SiON to the excimer laser lithography is proposed. It is exemplified by a process of fine gate processing with an SiON film restraining the reflectivity of a W (tungsten)--polycide film or an Al (aluminum) based material film.\nMeanwhile, after the patterning of the resist mask by such photolithography is finished, the antireflection film is etched in the subsequent etching process.\nIn this case, such a problem is now being apparent that the anisotropic shape of the underlying material layer may be deteriorated by oxygen discharged from SiON in the etching process, particularly in overetching. This problem is explained with reference to FIGS. 1 to 4.\nFIG. 1 shows a state of a wafer prior to the etching, in which a gate SiO.sub.x film 22, a W-polycide film 25 and an SiON antireflection film 26 are sequentially stacked on an Si substrate 21, with a resist mask 27 patterned in a predetermined shape being formed thereon. The W-polycide film 25 is composed of, from the bottom, a polysilicon layer 23 containing impurities and a WSi.sub.x (tungsten silicide) layer 24 which are sequentially stacked.\nIf the W-polycide film 25 is etched using a Cl.sub.2 /O.sub.2 mixed gas, the etching is promoted by a formation of etching reaction products, such as SiCl.sub.x and WClO.sub.x. On the other hand, a carbon based polymer derived from decomposition products provided by forward sputtering of the resist mask is deposited to form a sidewall protection film 28 on the sidewall surface of the pattern. If the wafer temperature is sufficiently low, SiCl.sub.x of relatively low vapor pressure among the etching decomposition products can be a component of the sidewall protection film 28.\nAs a result, a gate electrode 25a of anisotropic shape is formed at the end of just etching, as shown in FIG. 2. In FIG. 2, materials after the etching are denoted by their respective original numerals plus subscripts \"a\".\nHowever, if the overetching follows the just etching, regression of the edge of the resist mask 27 causes the SiON antireflection film 26a to have its end surface tapered to be easily exposed, as shown in FIG. 3. SiON, having an element composition ratio of approximately Si:O:N=2:1:1, is richer in Si than SiO.sub.2 is. Consequently, SiON has low durability to a Cl based plasma, and easily discharges active O* when its exposed end surface is etched. Then, O* removes the sidewall protection film 28 in the form of CO.sub.x, to lower sidewall protection effects. In addition, since the W-polycide film 25 to be etched is reduced in the overetching, a relatively excessive amount of O* is present in the etching gas.\nAs a result, a gate electrode 25b having an undercut is formed, as shown in FIG. 4. The material layers having the undercut denoted by their respective original numerals plus subscripts \"b\". The undercut is generated most conspicuously in the WSi.sub.x layer 24b. Since O* sputtered out from the end surface of the SiON antireflection film removes W atoms in the form of WClO.sub.x, the etchrate in the WSi.sub.x 24a is increased.\nAs the anisotropic shape of the gate electrode is thus deteriorated, serious problems rise, such as, the metallization resistance falling off the designed value and difficulty in forming the sidewall to attain an LDD structure.\nThe deterioration in the anisotropic shape in the overetching is not limited to the above-described SiON antireflection film, but is a phenomenon which may be generated in cases where an antireflection film capable of easily discharging oxygen is used and where conductive material layers of Al based metallization and the like other than the W-polycide film are used as etching targets."}
-{"text": "Multi-layered heteostructures are employed to implement devices for a number of applications. These applications include, but are not limited to, optoelectronic components (e.g. PIN junction or multi-quantum well). The functionality of these multi-layered heteostructures are typically built from layer to layer in a vertical direction, using different semiconductor materials. Further, the multi-layered heteostructures are vertically etched leading to the exposure of their sidewalls, and polymer is spun to seal the sidewalls. To facilitate provision of a contact to one of these devices, the polymer may be etched back to expose the top semiconductor layer, to allow a metal contact to be deposited thereon. Alternatively, a vertical via may be etched to open the polymer to facilitate contact between the top semiconductor layer and the metal contact.\nHowever, both practices have disadvantages. In particular, the former practice may not be able to clear the top semiconductor layer without exposing the sidewalls of some of the device layers underneath the top layer. Whereas, the latter practice is difficult and complicated, especially in the smaller than micro scale, e.g. at nanometer scale. As at the nanometer scale, not only alignment of the via mask becomes very difficult, making of the via mask in and of itself becomes almost impossible, due to current sub-micro lithography is unable to accurately resolve nanometer via printing. Also, at nanometer scale, the via approach will not allow the full use of the available area of the top semiconductor layer because a typical via approach requires some margin so the via must be smaller than the device. Even if the first practice is able to open the whole area of the top device layer, at micrometer or nanometer scale, the top semiconductor area may not be sufficiently large to provide a desired low contact resist interconnect (as resist is inversely proportional to the contact area)."}
-{"text": "1. Field of the Invention\nThis invention relates to a collapsible work horse having first and second pairs of legs pivotally mounted to a support beam to move from an extended or working position to a storage and/or transporting, e.g., collapsed, position, and a locking arrangement to lock the legs in the extended position and, more particularly, to a collapsible work horse having the legs secured in the extended position by a plunger mounted in each of the legs and biased into a hole in the support beam. The invention further relates to a work station having one or more work horses for supporting a shaping tool and for supporting the pieces to be shaped.\n2. Discussion of the Technical Problems\nIn general, work horses, also known as sawhorses or trestles, include a first pair of legs secured to one side of a support beam and a second pair of legs secured to an opposite side of the support beam. The legs can be fixedly secured to the support beam using fasteners, e.g. but not limited to, nails, screws, and/or nut and bolt arrangements, or detachably secured to the support beam using clamps. In general, the clamps include a pair of elongated members pivotally mounted together such that moving one end of the members away from one another moves the opposite ends of the members toward one another against the support beam. In another arrangement, the legs are secured by pivotally attaching the legs to the support beam as taught in U.S. Pat. No. 3,951,233 (hereinafter also referred to as \u201cU.S. Pat. No. '233\u201d).\nAlthough the presently available work horse designs are acceptable for their intended use, they have drawbacks. More particularly, work horses that have the legs and support beam fixedly secured together are usually moved and/or stored in the assembled state, which results in wasted unused space. The work horses that have the legs detachably secured to the support beam reduces the amount of unused space required for storage but requires disassembling the work horse, keeping track of the disassembled parts, and assembling the parts to use the work horse.\nThe collapsible work horse of U.S. Pat. No. '233 eliminates many of the problems discussed above; however, the work horse of U.S. Pat. No. '233 has limitations. More particularly, the extended legs of the work horse disclosed in U.S. Pat. No. '233 are maintained in the extended position by a constant frictional force applied to the pivot point of the legs. The frictional force is applied by tightening the bolt at the pivot point. For a detailed discussion of the arrangement to maintain the legs in the extended position, reference can be made to Patent '233.\nAs can be appreciated, tightening bolts to secure the legs in the extended position requires the use of the tool to tighten the bolts to secure the legs in the extended position and to loosen the bolts to move the legs to the collapsed position. It can be appreciated by those skilled in the art that it would be advantageous to provide a work horse that has legs that can be moved between the extended position and the collapsed position and does not have the drawbacks and/or limitations of the presently available work horses."}
-{"text": "1. Field of the Disclosure\nThe present disclosure relates to a bladed rotor, and more particularly relates to a bladed rotor for a turbo-machine such as a gas turbine engine. The disclosure is particularly suited for use in gas turbine compressor rotors, although it is to be appreciated that the disclosure is not limited to compressor rotors and could find application in other types of bladed rotors for use in other types of turbo-machines.\n2. Description of the Related Art\nConventional axial compressor rotors for gas turbine engines typically comprise a number of discs which are bolted or welded together to form an integral rotatable drum. Each disc can be considered to represent a central hub around which a plurality of rotor blades of aerofoil configuration are mounted. Each rotor blade is normally attached to the hub using a mechanical connection known as a root fixing. One such type of arrangement involves axially fixing the rotor blades to the periphery of the hub and involves the provision of a series of slots which are machined into the peripheral region of the hub and which are generally elongate parallel to one another. The slots are typically arranged so that they extend in a lengthwise direction which makes an acute angle of between 10 and 30 degrees to the rotational axis of the hub. Each slot is configured to receive a dove-tail or fir-tree shaped root fixing of a respective rotor blade.\nA radially outwardly biased sprung retaining ring is normally used to secure the root portions of the rotor blades within their respective mounting slots. The retention ring locates within radially inwardly open grooves formed around the hub at positions located between the blade mounting slots, under its radially outward bias. Similar grooves are provided on the rotor blades and so the retaining ring also locates in the blade grooves to axially retain the root portions of the blades in the mounting slots.\nIt is important for integrity reasons that during operation of the rotor that the retaining ring does not apply radial load to the blades within the blade grooves. The retaining ring must at all times remain radially inwardly spaced from the radially outmost region of each blade groove by a clearance gap. It is therefore normal to configure the arrangement such that the retaining ring only bears against the radially outmost regions of the hub grooves.\nHowever, it has been found that during service the retaining rings of the above-described type of axial fixing arrangement can be susceptible to wear on their radially outmost surfaces, as also can the inner surfaces of the hub grooves within which the rings locate. Over time, this wear can reduce the size of the radial clearance gap between the retaining ring and the blade grooves which, as indicated above, cannot be allowed to occur due to integrity concerns."}
-{"text": "The invention relates to a weft stop motion, or detector, for looms in which the sensor element responds or is sensitive to an electric charging of the weft thread without contact and also to looms in which the weft break stop motion of the present invention is used.\nA concept for thread detection is known from German Patent Specification 3,758,403, for example. Various embodiments of electrostatic transformers are disclosed therein. These sensors are mainly used in air-jet looms. The weft thread is electrically charged during its removal from the weft thread supply because of the resultant friction and also during the weft insertion because of friction with the air. The electrostatic detection registers the presence of a textile fiber which is moving past and is electrically charged in this way, and the passage of the tip of an inserted weft thread in particular can also be detected. Weft break stop motions are used in the weft channel of the loom. The known embodiments are relatively large and heavy and are frequently constructed in the form of a confusor drop wire. The high-speed air-jet looms having a correspondingly high beat-up speed of the reed which are commonly used nowadays produce high vibration and acceleration loads on the known weft break stop motions, so that the known embodiments are no longer suitable for use on air-jet looms or the resultant electrical signals are very noisy."}
-{"text": "Many disabled and elderly persons are inhibited or precluded from playing golf due to various aspects of their conditions. Some may have trouble getting into position to swing a golf club, such as those who rely on a wheelchair for mobility. Others may not be able to effectively grip and/or move their body to swing the club, such as persons suffering from certain paralysis. Still others may be physically able, but lack the cognitive ability to swing a golf club in a traditional manner, such as some mentally disabled or autistic persons.\nAdditionally, many persons seek to improve their golf swing through the use of training apparatuses. Training apparatuses physically manipulate a person's movement or the movement of their club to teach certain swing mechanics. Many persons learn more effectively by witnessing visual demonstrations of certain techniques including, but not limited to, the pendulum-like swing motion often used in chipping and putting."}
-{"text": "Vehicle mounted cable reel handling apparatus adapted to be carried on a trailer vehicle are shown in U.S. Pat. Nos. 3,091,413 and 3,063,584. Reel handling apparatus has also been provided for pickup trucks as disclosed in U.S. Pat. Nos. 3,165,214; 3,184,082; 3,036,790 and 3,325,118. In the above patents, the apparatus generally includes vertically movable lift arms pivotally connected to the vehicle for engaging and transferring a single ground supported reel onto the vehicle for transport. The lift arms are engageable with the reel at all times even when the reel is in the transport position.\nIn some of the reel handling structures for trucks to handle a pair of reels for transport, the operating cylinders for the reel lift arms are arranged on the truck bed so as to appreciably limit the space for reel storage as appears in U.S. Pat. Nos. 2,876,916 and 3,902,612. The Anderson U.S. Pat. No. 3,625,380 and McVaugh U.S. Pat. No. 3,820,673 use front and rear pairs of lift arms, with a first reel lifted from the ground by rear arms being transferred to the front lift arms and carried thereon to a farward transport position. The second reel when lifted from the ground remains on the rear lift arms for transport in a position adjacent to and rearwardly of the first reel.\nAlthough transfer of a reel from the rear lift arms to the front lift arms was generally satisfactory, the double lift arm arrangement was relatively expensive and difficult to accommodate within the limited space requirements on the truck bed, especially as restricted with the growing demand for larger side mounted tool carrying compartment units. With the compartment units extended from the truck cab to positions over and behind the truck rear wheel and axle assembly space requirements for transporting the reels become more critical.\nThe Hall U.S. Pat. No. 3,902,612 partially solves this problem by using transversely spaced tiltable beams extended longitudinally of the truck for receiving a reel from a pair of rear lift arms. On a controlled downward and forward tilting movement of the beams the transferred reel is rolled by gravity action to a forward transport position. However, by virtue of the lift arms being actuated by cylinders mounted on the truck bed, the transverse distance between the beams is appreciably reduced. As a result the reel has a spindle of reduced length, when lifted from the ground, which is then replaced by a longer spindle before the reel can be supported on the beams. Hall, therefore, has no provision for side compartment units and requires a manual changing of reel spindles, and a manual actuation of the tiltable beams to roll a reel to a front transport position. These disadvantages of the Hall apparatus are eliminated by the apparatus of this invention."}
-{"text": "1. Field of the Invention\nThe present invention relates to a digital camera having a self-timer shooting function and an image display function.\n2. Description of the Related Art\nIn resent years, digital cameras have rapidly been becoming widespread, and various types of digital cameras are supplied to the market. Generally, a digital camera frequently has an image display device comprising a liquid crystal display (hereinafter, referred to as an LCD) or the like on its back. Because of the LCD, the photographer can confirm the image to be shot without viewing through the finder. Therefore, it is possible to shoot the subject at an extremely free angle. Moreover, a conventional digital camera is provided with a function to continue displaying the shot image on the image display device for a predetermined time every time one frame is shot, that is, a function to hold the shot image for a predetermined time. Because of this function, the photographer can check whether the shot image is desired or not without performing any complicated operations.\nWhether silver halide film cameras or digital cameras, cameras are frequently provided with a function to perform shooting by use of a self-timer (hereinafter, referred to as a self-timer function). According to the self-timer function which is used, for example, when the photographer himself or herself is the subject to be shot, shooting is on standby for a predetermined time after the depression of the release button, and shooting is performed after the predetermined time has elapsed. Because of this function, the photographer can shoot himself or herself by moving to the position of shooting within the predetermined time after depressing the release button.\nHowever, since the photographer is the subject when self-timer shooting is performed as mentioned above, it takes time for the photographer to return to the camera after shooting is performed. Therefore, even though the digital camera has the image display device and the function to display the shot image for a predetermined time, it frequently occurs that the display of the shot image ends before the photographer returns to the camera to check the display device on the back of the camera. That is, the photographer cannot check the shot image. Since recorded images can be read out from a recording medium and played back, the shot image can be checked by playing it back. However, complicated operations-such as switching to a reproduction mode and specification of the frame to be played back are necessarily performed every time, which is inconvenient.\nIn view of the above-mentioned problem, an object of the present invention is to provide a digital camera in which the photographer can easily check whether the shot image is desired or not without performing any complicated operations even when self-timer shooting is performed."}
-{"text": "1. Field of the Invention\nThis invention relates to a suction cup device, more particularly to a suction cup device which defines a volume-variable space to produce an adjustable reduced pressure with a desired suction strength.\n2. Description of the Related Art\nFIG. 1 shows a conventional suction cup device 10 disclosed in U.S. Patent Publication No. US 2010/0116954 A1, which is adapted to be attached to a flat wall 1, and which includes a suction cup 13 with a rack 14 connected to the center thereof, a mount body 11 attached to the suction cup 13 by virtue of a pressing ring 12, and a lever 15 pivotably mounted on the mount body 11 and provided with a pinion 152 meshing with the rack 14. By turning the lever 15, a central part of the suction cup 13 can be pulled away from the flat wall 1 to produce a reduced pressure in an enclosed space between the flat wall 1 and the suction cup 13 for holding the suction cup device 10 against the flat wall 1. However, in the conventional suction cup device 10, the strength of suction force generated as a result of movement of the central part of the suction cup 13 is constant, and cannot be increased after a period of use, thereby resulting in undesired disengagement of the suction cup device 10 from the flat wall 1.\nAnother conventional suction cup device is disclosed in U.S. Pat. No. 7,021,593 B1, which includes a suction lock/release disposed on a lever and deep scoops disposed on the mount body to lock the lever at a fixed angle so as to adjust the suction strength of a suction cup. However, assembly of the suction lock/release to the lever is troublesome. Besides, manual operation of the suction lock/release is required to lock or unlock the lever, thereby rendering the adjustment complicated."}
-{"text": "1. Field of the Invention\nThe invention relates to a process for measuring the flow vectors in gas currents containing optically detectable particles wherein a focusing means focuses at least two spatially separated beams in at least two focusing point in a measuring volume.\n2. Description of the Related Art\nProcesses for measuring flow vectors in gas currents have been known to employ light of a light source focused by a focusing means in the flow channel at two focusing points positioned in a close succession (U.S. Pat. No. 3,941,477). Particles contained in the gas current are illuminated in traversing the focusing points. Due to the stray radiation reflected by the particles, a start pulse is produced when the first focusing point is traversed, while a stop pulse is produced during the traversing of the second focusing point. From the time interval between the start pulse and the stop pulse, it is possible to determine the component of the particle speed vector in direction of the straight line traversing the focusing points. By moving the focusing device, said direction may be varied thus permitting the detection of flow vectors having different directions. However, by said method, it is only possible to measure the components of the flow vectors which extend in a normal plane relative to the optical axis of the focusing means. The component extending in the direction of the optical axis may not be determined. Said process is designated to 2d-process to refer to the two-dimensional vector measurement.\nA further development of said process is designated as 3d-process, by which the vector component extending in direction of the optical axis may be detected as well (British Pat. No. 2,109,548) and in which four laser beams are used two of which each form a beam pair. The beams of each pair being directed to two focusing points situated in the same normal plane relative to the optical axis. Due to the differences of the direction of incidence of the beams directed to a focusing point, the flow directions measured by means of the beam pairs are determined differently. From said difference of direction, one may draw a conclusion concerning the flow component in direction of the optical axis of the system. The expenditure and the laser capacity required by said known process are quite considerable."}
-{"text": "Many processes in biology, including transcription, translation, and metabolic or signal transduction pathways, are mediated by noN-covalently-associated multienzyme complexes1, 101. The formation of multiprotein or protein-nucleic acid complexes produce the most efficient chemical machinery. Much of modern biological research is concerned with identifying proteins involved in cellular processes, determining their functions and how, when, and where they interact with other proteins involved in specific pathways. Further, with rapid advances in genome sequencing projects there is a need to develop strategies to define \u201cprotein linkage maps\u201d, detailed inventories of protein interactions that make up functional assemblies of proteins2,3. Despite the importance of understanding protein assembly in biological processes, there are few convenient methods for studying protein-protein interactions in vivo4,5. Approaches include the use of chemical crosslinking reagents and resonance energy transfer between dye-coupled proteins102, 103. A powerful and commonly used strategy, the yeast two-hybrid system, is used to identify novel protein-protein interactions and to examine the amino acid determinants of specific protein interactions4,6-8. The approach allows for rapid screening of a large number of clones, including cDNA libraries. Limitations of this technique include the fact that the interaction must occur in a specific context (the nucleus of S. cerevisiae), and generally cannot be used to distinguish induced versus constitutive interactions.\nRecently, a novel strategy for detecting protein-protein interactions has been demonstrated by Johnsson and Varshayskyl108 called the ubiquitin-based split protein sensor (USPS)9. The strategy is based on cleavage of proteins with N-terminal fusions to ubiquitin by cytosolic proteases (ubiquitinases) that recognize its tertiary structure. The strategy depends on the reassembly of the tertiary structure of the protein ubiquitin from complementary N- and C-terminal fragments and crucially, on the augmentation of this reassembly by oligomerization domains fused to these fragments. Reassembly is detected as specific proteolysis of the assembled product by cytosolic proteases (ubiquitinases). The authors demonstrated that a fusion of a reporter protein-ubiquitin C-terminal fragment could also be cleaved by ubiquitinases, but only if co-expressed with an N-terminal fragment of ubiquitin that was complementary to the C-terminal fragment. The reconstitution of observable ubiquitinase activity only occurred if the N- and C-terminal fragments were bound through GCN4 leucine zippers109,110. The authors suggested that this \u201csplit-gene\u201d strategy could be used as an in vivo assay of protein-protein interactions and analysis of protein assembly kinetics in cells. Unfortunately, this strategy requires additional cellular factors (in this case ubiquitinases) and the detection method does not lend itself to high-throughput screening of cDNA libraries.\nRossi, F., C. A. Charlton, and H. M. Blau (1997) Proc. Nat. Acad. Sci. (USA) 94, 8405-8410) have reported an assay based on the classical complementation of \u03b1 and \u03c9 fragments of \u03b2-galactosidase (\u03b2-gal) and induction of complementation by induced oligomerization of the proteins FKBP12 and the mammalian target of rapamycin by rapamycin in transfected C2C12 myoblast cell lines. Reconstitution of b-gal activity is detected using substrate fluorescein di-\u03b2-D-galactopyranoside using several fluorescence detection assays. While this assay bears some resemblance to the present invention, there are several significant distinguishing differences. First, this particular complementation approach has been used for over thirty years in a vast number of applications including the detection of protein-protein interactions. Krevolin, M. and D. Kates (1993) U.S. Pat. No. 5,362,625) teaches the use of this complementation to detect protein-protein interactions. Also achievement of \u03b2-gal complementation in mammalian cells has previously been reported (Moosmann, P. and S. Rusconi (1996) Nucl. Acids Res. 24, 1171-1172). The individual PCAs presented here are completely de novo designed interaction detection assays, not described in any way previously except for publications arising from applicants laboratory. Secondly, this application describes a general strategy to develop molecular interaction assays from a large number of enzyme or protein detectors, all de novo designed assays, whereas the \u03b2-gal assay is not novel, nor are any general strategies or advancements over previously well documented applications given.\nAs in the USPS, the yeast-two hybrid strategy requires additional cellular machinery for detection that exist only in specific cellular compartments. There is therefore a need for a detection system which uses the reconstitution of a specific enzyme activity from fragments as the assay itself, without the requirement for other proteins for the detection of the activity. Preferably, the assay would involve an oligomerization-assisted complementation of fragments of monomeric or multimeric enzymes that require no other proteins for the detection of their activity. Furthermore, if the structure of an enzyme were known it would be possible to design fragments of the enzyme to ensure that the reassembled fragments would be active and to introduce mutations to alter the stringency of detection of reassembly. However, knowledge of structure is not a prerequisite to the design of complementing fragments, as will be explained below. The flexibility allowed in the design of such an approach would make it applicable to situations where other detection systems may not be suitable.\nRecent advances in human genomics research has led to rapid progress in the identification of novel genes. In applications to biological and pharmaceutical research, there is now the pressing need to determine the functions of novel gene products; for example, for genes shown to be involved in disease phenotypes. It is in addressing questions of function where genomics-based pharmaceutical research becomes bogged down and there is now the need for advances in the development of simple and automatable functional assays. A first step in defining the function of a novel gene is to determine its interactions with other gene products in an appropriate context; that is, since proteins make specific interactions with other proteins or other biopolymers as part of functional assemblies, an appropriate way to examine the function of a novel gene is to determine its physical relationships with the products of other genes.\nScreening techniques for protein interactions, such as the yeast \u201ctwo-hybrid\u201d system, have transformed molecular biology, but can only be used to study specific types of constitutively interacting proteins or interactions of proteins with other molecules, in narrowly defined cellular and compartmental contexts and require a complex cellular machinery (transcription) to work. To rationally screen for protein interactions within the context of a specific problem requires more flexible approaches. Specifically, assays that meet criteria necessary not only to detecting molecular interactions, but also to validating these interactions as specific and biologically relevant.\nA list of assay characteristics that meet such criteria are as follows:\n1) Allow for the detection of protein-protein, protein-DNA/RNA or protein-drug interactions in vivo or in vitro.\n2) Allow for the detection of these interactions in appropriate contexts, such as within a specific organism, cell type, cellular compartment, or organelle.\n3) Allow for the detection of induced versus constitutive protein-protein interactions (such as by a cell growth or inhibitory factor).\n4) To be able to distinguish specific versus non-specific protein-protein interactions by controlling the sensitivity of the assay.\n5) Allow for the detection of the kinetics of protein assembly in cells.\n6) Allow for screening of cDNA, small organic molecule, or DNA or RNA libraries for molecular interactions."}
-{"text": "1. Technical Field\nThe present invention is directed generally to wireless communication systems and, more particularly, to a system and method for parameter selection to avoid interference in a wireless communication system.\n2. Description of the Related Art\nEarly wireless communication devices, commonly known as cell phones, provided wireless voice services to the user. These early phones have been replaced with wireless communication devices capable of delivering voice, data, and multi-media information. In addition, wireless devices often include location determination using the Global Positioning System (GPS). The delivery of these additional services requires additional bandwidth. In some cases, bandwidth previously allocated for one purpose has been reassigned for the delivery of wireless communication services. For example, the spectrum originally allocated to Ultra-High Frequency (UHF) television has been partially reallocated for wireless communication services.\nDevices are being designed with multiple services that depend on multiple radio systems being operated at the same time. For example, devices are being designed that can connect to the cellular network using several different radio protocols and frequency bands. In addition, these devices may have other applications, such as broadcast television or Bluetooth, which use independent radio systems.\nThese independent radio systems may interfere with, or be interfered by, the radio system used for cellular operation. One can appreciate that the operation of multiple transceivers within a single device may decrease the operational capability of the device. Therefore, it can be appreciated that there is a significant need to reduce interference among the multiple transceiver systems. The present invention provides this, and other advantages, as will be apparent from the following detailed description and accompanying figures."}
-{"text": "The present invention relates to a circular comb for combing machines with a segment-shaped basic element on which the individual needles (teeth) are located parallel to each other and are connected by pressure strips fastened to the basic element under pressure.\nIn modern combing technology, the teeth (needles) are no longer fastened to a barrette (or bar), but to a wire or strip-like needle carrier to which they are soldered, welded or glued. For fastening such needle strips to a circular comb in a known embodiment (e.g. German Patent DT-OS No. 2,002,020), wedge-shaped grooves are located on the outside circumference of the basic element at predetermined intervals. Needle strips are inserted into these grooves and are pressed by means of a wedge-shaped clamping strip against the webs of the basic element remaining between the grooves. The necessary pressure is achieved by a number of screws distributed in the lengthwise direction over the clamping strip. These screws are threaded into the basic element. Since the combs are relatively close together and the needle tips in the peripheral direction are only 8 mm apart at the outside circumference of the basic element, and since the clamping strips are tapering conically inward, only very small screws can be used for fasteners. The same applies with respect to the pressure screws for detaching the clamping strips from the basic element; additional threads for these must be located in the clamping strip. When replacing the needle strips, a large number of small screws must be unscrewed and tightened again; this requires a considerable expenditure of time.\nIt is also known in the art how to screw fasten the individual needle strips one after the other to the basic element of a circular comb. The first needle strip is placed on the outside of the basic element and fastened by means of a wedge-shaped strip through which the screws pass to engage the threads in the basic element. The needle strip is pressed between the basic element and the wedge-shaped strip. This procedure is repeated till the last needle strip of the segment is reached; it is followed by a final segment which is designed so that the circular comb can be mounted on the machine shaft. The manufacture of such a circular comb is expensive. Even greater is the disadvantage that the individual needle strips cannot be exchanged one for one. In the extreme case, all preceding needle strips of such a circular comb must be removed before the last needle strip can be detached and replaced. This results in extremely cumbersome shutdowns.\nIt is, therefore, an object of the present invention to simplify the fastening of the needle strip on the circular comb of a combing machine and to provide individual interchangeability of the individual needle strips, regardless of whether it is the first or the last or any other needle strip.\nAnother object of the present invention is to provide a circular comb arrangement which may be economically fabricated and maintained in service.\nA further object of the present invention is to provide a circular comb arrangement, as described, which has a substantially long operating life."}
-{"text": "1. Field of the Invention\nThe present invention relates to an errorproof device, and more particularly to the combination of an errorproof device and a modular socket.\n2. Description of Related Art\nA conventional communication modular socket for connection with a modular plug used in a telephone line or a modem does not have the ability to distinguish whether the plug to be inserted into the socket has the appropriate dimension. For example, the currently available RJ 11 or RJ45 plugs are both used for communication devices and respectively have a dimension different from the other. Because the RJ 11 plug has a smaller dimension than that of the RJ 45 plug, the RJ 11 plug may be erroneously inserted into the modular socket (namely the RJ 45 socket) configured to mate for the RJ 45 plug and thus leads on the RJ 45 modular socket may be damaged.\nTo overcome the shortcomings, the present invention tends to provide an improved modular socket having an errorproof device therein to mitigate the aforementioned problems."}
-{"text": "There exist in the marketplace today a number of different hook-fastener media to be described below. It is our belief that each of these existing hook-fasteners suffers from one or more shortcomings which hamper their utility and utilization."}
-{"text": "1. Field of the Invention\nThe present invention is directed generally to a cosmetic application system and method, and more particularly to a system and method for removing excess mascara from a mascara brush upon withdrawal from a container.\n2. Description of the Related Art\nVarious techniques and structures have been used to reduce the amount of mascara on a mascara brush upon removal from a container. However, a number of disadvantages associated with these techniques and structures has inhibited their widespread use and manufacture.\nIn particular, U.S. Pat. Nos. 4,194,848, 4,332,494, 4,407,311, 4,609,300 and 4,705,053 are directed to mascara applicators having a complex structure for varying the amount of mascara remaining on a brush after removal from a container. A flexible member is disposed in the neck of the container to provide some degree of variation in the amount of mascara removed from a brush as it passes through an opening in the container. However, each of these patents is directed to a complex structure, which is difficult and costly to manufacture. Moreover, many of these structures do not facilitate continuous variation of the amount of mascara to be removed from a brush. In addition, because these structures apply an equal force against the brush during removal and re-insertion of the brush into the container, these systems unnecessarily impede a user's ability to reinsert the brush into the container after each use. In U.S. Pat. No. 5,397,193, an additional attempt was made to provide a system for removing excess mascara from an applicator brush. In this system, a plurality of internal flexible bristles are used to remove excess mascara from the applicator brush. As with the aforementioned patents, this system is also costly and difficult to manufacture, and does not facilitate continuous variation in the amount of force to be applied to the mascara brush upon removal from its container. In addition, this system also unnecessarily impedes a user's ability to reinsert the brush into the container after each use.\nIn addition, each of the aforementioned systems, because of their complicated internal structure, is particularly difficult to clean. Accordingly, these go systems do not lend themselves for use with any form of reusable or interchangeable mascara system."}
-{"text": "1. Field\nThe invention relates to the field of computer configuration and, more particularly, to the storage of computer configuration signals.\n2. Background Information\nComputer systems may store configuration signals in a memory. A computer system is any device including a processor capable of executing one or more instructions to generate signals. Such signals typically take the form of sequences of binary signals known as bits. Examples of computer systems are personal computers, workstation computers, server computers, hand held computers, and set top boxes to name just a few examples. Configuration signals are signals that may determine various settings for the operation of the computer system. For example, configuration signals may determine whether various input/output (I/O) ports comprised by the system are enabled, and I/O addresses for these ports. Configuration signals may determine other computer system settings as well. Such computer configuration signals are well known in the art as \u201cset up information\u201d. On personal computers, setup information is also often stored in a memory known as a real time clock (RTC) complementary metal oxide silicon (CMOS).\nSetup information may be applied prior to or during the booting of a program to control the computer system. Booting is the process of placing a sequence of instructions (a program) in control of various computer system resources. Resources include memory, interrupts, files, and I/O ports. An example of a program to boot is an operating system. An operating system is a program which controls various computer resources including those mentioned previously and further including typical I/O devices, such as a mouse and keyboard. Examples of operating systems are the Unix\u2122 operating system and the Microsoft\u2122 Windows\u2122 operating system.\nSetup information may be read, altered, and written back to a CMOS or other memory, where it is stored using a special program called a \u201csetup program.\u201d The setup program may be part of the sequence of instructions comprising the computer system's power-on self test (POST) program. Often, the POST executes prior to the basic input/output system (BIOS) program of the computer system in order to initialize settings.\nThe settings determined by setup information may vary among different computer makes and models. Furthermore, the location and lengths of the bit sequences that comprise the setup information in the memory in which they are stored may vary. Accordingly, it may be difficult to create one set up program to read, alter, and write back set up) information for various makes and models of computer systems. Instead, multiple set up programs may be called for different makes and models of computer system.\nExisting setup programs typically employ a crude \u201ctextual\u201d interface. Textual interfaces are well known in the art and may comprise an 80\u00d725 matrix of character positions. The number, type, and position of characters in a textual interface is limited as are the number of colors in which such characters may be displayed. It is well known that such textual interfaces are more limited than modem \u201cgraphical user interfaces\u201d (GUI) which provide individual control of the color and position on a per pixel basis on the computer system display. Typically, it is the operating system which implements a graphical user interface for the computer system. However, setup programs may execute before the operating system boots, and, therefore, the setup program may employ a less sophisticated textual interface instead of a GUI.\nThus, there is a continuing need for a setup program which may operate with various makes and models of computer systems and which may take advantage of graphical user interface features provided by an operating system."}
-{"text": "1. Field of the Invention\nThis invention relates to the field of ceramics and particularly to ZrO.sub.2 ceramics.\n2. Description of the Prior Art\nDuring cooling, ZrO.sub.2 undergoes a martensitic-type transformation from a tetragonal crystal structure to a monoclinic crystal structure with a concurrent increase in volume and an anisotropic shape change. For pure ZrO.sub.2 the transformation begins at about 1200.degree. C. and proceeds until complete at about 600.degree. C.\nAttempts have been made to utilize this transformation in order to improve the fracture toughness of ceramic composites. In one approach, ZrO.sub.2 particles have been added to an Al.sub.2 O.sub.3 matrix to form a second phase dispersion (N. Claussen, J. Am. Ceram. Soc. 59, pg. 49, 1976). Expansion and shape change of the ZrO.sub.2 as it transformed from the high temperature tetragonal structure to the room temperature monoclinic structure created microcracks. The resulting increase in fracture toughness was attributed to energy absorption by these microcracks.\nMore recently, attempts have been made to increase the toughness of ZrO.sub.2 ceramics by taking advantage of metastable grains of tetragonal ZrO.sub.2 within a surrounding matrix. These are grains of ZrO.sub.2 which are tetragonal rather than monoclinic despite the fact that their temperature is below the unconstrained equilibrium transformation temperature range.\nThe metastable condition can be obtained by surrounding fine grains of ZrO.sub.2 in a constraining matrix such as Al.sub.2 O.sub.3. The matrix constrains the volume and shape change associated with the transformation of the grains of ZrO.sub.2 and holds the ZrO.sub.2 in its tegragonal state.\nThe tetragonal grains of ZrO.sub.2 increase the fracture toughness of the ceramic composite by increasing the energy required for a crack to propogate. If a crack starts in the ceramic composite, the metastable grains of tetragonal ZrO.sub.2 in the stress field adjacent the crack transform to the stable monoclinic structure. The work done by the applied stresses to reduce this transformation is loss and thus the stress-induced transformation increases the material's fracture toughness.\nMetastable tetragonal grains of ZrO.sub.2 have been observed in an Al.sub.2 O.sub.3 /ZrO.sub.2 ceramic composite containing 17 volume % ZrO.sub.2 (N. Claussen, J. Am. Ceram. Soc. 59, pg. 85, 1978). However, to maintain the metastable tetragonal structure, the ZrO.sub.2 grains had to be less than about 0.5 .mu.m in diameter. Larger grains transformed to the stable monoclinic structure. Additional work has shown that the amount of metastable tetragonal ZrO.sub.2 that can be retained in the matrix decreases as the volume % of ZrO.sub.2 in the Al.sub.2 O.sub.3 /ZrO.sub.2 ceramic composite increases. Very little of the ZrO.sub.2 can be retained in the metastable tetragonal structure in Al.sub.2 O.sub.3 /ZrO.sub.2 composites having more than 20 volume % ZrO.sub.2. Such limitations of grain size and volume % of ZrO.sub.2 reduces the practicality and the toughness of prior art Al.sub.2 O.sub.3 /ZrO.sub.2 ceramic composites."}
-{"text": "Technical Field\nThe present disclosure relates to memory devices such as a semiconductor memory, and method for testing reliability of memory devices.\nRelated Art\nIn general reliability tests on semiconductor memories, a testing device is used to write and read data to and from all regions in a memory array with a known test pattern, and the data written to the memory array by the testing device (expected value) is compared with the data read from memory array by the testing device, so as to check the reliability of the memory array.\nIn pre-shipment inspection of semiconductor memories, in order to reduce a testing cost, reliability test is generally performed concurrently on multiple semiconductor memories, by connecting multiple semiconductor memories to one testing device and writing and reading data to and from the multiple semiconductor memories with a common test pattern.\nWith semiconductor memories provided with a pseudo-random number generator for improving security, random number values of pseudo-random numbers are predictable, and thus multiple semiconductor memories can generate an identical pseudo-random number by using a common algorithm. Pre-shipment inspection can therefore be conducted concurrently on multiple semiconductor memories with one testing device, in the same way as on general semiconductor memories.\nRandom number generators are cryptographic technology employed for a wide variety of uses in many security systems.\nRandom numbers generated by random number generators are used for, for example, key information in a cryptographic algorithm, or authentication codes for mutual authentication between devices, and are closely related to the security strength of a system and thus highly confidential information.\nRandom numbers generated by random number generators therefore need to be highly random. At shipment of semiconductor devices provided with a random number generator, a random number test is normally performed to evaluate whether a random number generator generates random numbers that meet a required level.\nJP2005-517998A and WO2005/124537A describe a technique to evaluate whether the frequencies of appearance of \u201c0\u201d and \u201c1\u201d in random numbers generated by a random number generator are within an allowable range."}
-{"text": "1. Field of the Invention\nThe invention relates to a combined power station installation with a gas turbine and a steam turbine, in which the exhaust gases from the gas turbine give up their residual heat to the steam turbine via the working medium flowing in a waste heat boiler, whereby the waste heat boiler consists essentially of an economizer, an evaporator and a superheater and whereby at least one cooling air cooler is provided which is designed as a forced circulation steam generator and is connected on the water side to the economizer of the waste heat boiler.\n2. Discussion of Background\nGas turbines of the modern generation and the higher power class operate with very high turbine inlet temperatures, which makes cooling of the combustion chambers, the rotors and the blading unavoidable. For this purpose, highly compressed air is generally extracted at the compressor outlet and, if appropriate, from a lower pressure stage. Because a very high proportion of the compressed air is consumed for the currently conventional premixed combustion, there remains--on the one hand--only a minimum of cooling air for cooling purposes. On the other hand, this air intended for cooling is already very hot because of the compression so that it is desirable that it should be previously cooled. Cooling by means of water injection (\"gas quenching\") is known for this purpose; in this method, however, the high-quality heat of the cooling air, whose proportion can amount to as much as 20 MW in current machines, is only partially utilized. In consequence, the use of forced circulation steam generators as coolers for recooling seems appropriate, particularly if the gas turbine operates in a combined gas/steam turbine process with waste heat steam generation.\nSuch a once-through steam generator for cooling highly compressed air of the type mentioned at the beginning is known, in association with a combined gas/steam turbine process, from EP-A-709 561. In this specification, a partial flow of the boiler feed water is extracted either upstream or downstream of the economizer and, after further preheating, evaporation and superheating in the cooler, is fed back into the high pressure superheater of the waste heat boiler. This boiler is designed as a circulating system boiler with drums. In order to avoid the penetration of moisture or water into the steam turbine when the cooler is run wet, the heated water or wet steam is fed into a blow-down tank until the cooler is dry or until defined conditions are stably present at the cooler outlet, for example hot steam with a few degrees Kelvin superheat or wet steam with a humidity of a few percent. In addition to the water losses, this has the consequent disadvantage of a corresponding monitoring and control system."}
-{"text": "In a time-of-flight (TOF) range sensor configured to acquire a range image, by using TOF scheme, a potential just under a gate electrode of a MOS structure is controlled in a vertical direction (depth direction) of the MOS structure. For example, a CMOS distance-measuring element and a TOF image sensor using the CMOS distance-measuring element are disclosed in patent literature (PTL) 1. The CMOS distance-measuring element has a structure such that an n-type charge-generation buried region buried in a p-type semiconductor layer, a charge-transfer buried region, a charge-read-out buried region, an insulating film covering the charge-generation buried region and the charge-transfer buried region, a transfer gate electrode provided on the insulating film to transfer signal charges to the charge-transfer buried region, a read-out gate electrode provided on the insulating film to transfer the signal charges to the charge-read-out buried region. When pulse lights are irradiated to the charge-generation buried region in the CMOS distance-measuring element recited in PTL 1, light signals are converted to signal charges in a semiconductor layer just under the charge-generation buried region, and a distance to an object is measured from a distribution ratio of charges accumulated in the charge-transfer buried region.\nThe CMOS distance-measuring element or the TOF image sensor using the CMOS distance-measuring element has a problem of generation of noise or dark current caused by interface defects, interface states, or the like just under the transfer gate electrode. In addition, in the case of using the transfer gate electrode disclosed in PTL 1, a potential gradient over a long distance is difficult to control, and a substantially uniform electric field over a long distance of a charge transport path is practically difficult to maintain. For this reason, in the charge-modulation elements such as distance-measuring elements having long charge transport paths, carriers are stopped in the middle of the charge transport paths, and thus, there are issues such that expected performances of the charge-modulation element are difficult to achieve."}
-{"text": "1. Field of the Invention\nThe present invention relates to an automatic answering method and apparatus for supporting a question reply process of replying to a question document of a text format.\n2. Description of the Related Art\nWith recent widespread of computerization, questions to companies or the like are often made by form inputs at home pages or e-mails. If every question is to be answered manually on the company side, many operators are required and the cost increases. A novice operator can not answer some questions or it takes a long time for the novice operator to answer a question. In order to solve this problem, a question-answering system has been introduced recently. With this system, a question document is input and its content is analyzed to select a reply example candidate from reply examples and question-reply examples prepared for each question content and to present the selected reply example candidate.\nMost of such question-answering systems assume, however, that one document contains only one consultation content. Therefore, if a plurality of question contents are written in one document, the systems cannot analyze each question content, resulting in a low reply precision.\nAnother technique is disclosed in JP-A-2002-132661. This technique discloses means for dividing one document containing a plurality of question contents, into each question content. The divided question content is analyzed to select a reply example candidate. A reply precision representative of a likelihood or degree of each reply example candidate for the question content is calculated. If the reply precision has a predetermined value or higher, an answer is formed from the reply example candidate, whereas if the reply precision is lower than the predetermined value, an instruction is given to compose a new answer.\nThe conventional technique disclosed in JP-A-2002-132661 describes that the means for dividing a document into each question content performs a division process by using \u201cnumber\u201d, \u201calphabet\u201d, \u201c.\u201d, an indent, a conjunction such as \u201cor\u201d, and the like. However, if a document is divided into each question content by using \u201cnumber\u201d, an indent and the like as a separator, there occurs the problem that one question content is divided into a plurality of sentences. Conversely, there arises the problem that if the range of a question content is broad, example candidates for a plurality of question contents cannot be selected.\nAccording to conventional techniques, since a question document is divided basing upon only the information about the contents of the question document, the divided range may not be covered by each reply example candidate. Namely, it is necessary to divide a question document so as to be covered by a prepared reply example candidate, and not to divide it by referring only to the question document content.\nSince a question document divided basing upon conventional techniques may be a document irrelevant to the question document content, the reply example candidate generation process is adversely affected so that the reply example candidate generation precision lowers. It also takes a time for a reply composition operator to find a proper document to be read.\nAccording to conventional techniques, a reply precision representative of the likelihood value of a reply example candidate is calculated, and if the reply precision is a predetermined value or higher, a reply is generated from the reply example candidate to automatically answer (automatically return) the question. If the reply precision is lower than the predetermined value, an instruction is given to compose a new answer. However, if there are a large number of types of replies or if a similar question requires a different answer, the reply precision lowers so that the number of samples exceeding a predetermined threshold reduces. Therefore, the number of samples capable of being used for the automatic reply reduces, and the number of cases requiring to generate new answers increases. There arises the problem of a low operator work efficiency or an automatic reply using an erroneous reply example candidate."}
-{"text": "This invention relates to a fluid drive oil recovery process which utilizes an injection of CO.sub.2, surfactant and water into subterranean oil reservoir. More particularly, the invention relates to such a process in which the surfactant selected for use is a particular member of a relatively highly chemically stable and salt-tolerant class of surfactants and is uniquely suited for use in the reservoir to be treated.\nNumerous patents have been issued on materials and techniques which are pertinent to an oil recovery process that utilizes an injection of CO.sub.2, surfactant and water. The U.S. Pat. Nos. 2,226,119; 2,233,381 and 2,233,382 describe polyalkoxylated alcoholic or phenolic surfactants which are generally useful in aqueous liquid fluid drive oil recovery processes. U.S. Pat. No. 2,623,596 indicates that an increased oil recovery may be obtained by a fluid drive process which injects highly pressurized CO.sub.2. U.S. Pat. No. 3,065,790 indicates that, in a fluid drive process, the cost effectiveness of highly pressurized CO.sub.2 may be increased by injecting a slug of the CO.sub.2 ahead of a cheaper drive fluid. U.S. Pat. No. 3,330,346 indicates that almost any process for forming foam within a reservoir may be improved by using as the surfactant a polyalkoxylated alcohol sulfate of an alcohol containing 10 to 16 carbon atoms. U.S. Pat. No. 3,342,256 indicates that, in a fluid drive process, the oil-displacing efficiency of a CO.sub.2 slug may be increased by including water and a foaming surfactant within that slug. U.S. Pat. No. 3,529,668 indicates that, in a fluid drive process, the efficiency of a slug of foamed CO.sub.2 may be increased by displacing it with specifically proportioned alternating slugs of gas and liquid. U.S. Pat. No. 4,088,190 indicates that, in a fluid drive process, the heat stability and durability of a CO.sub.2 foam may be increased by using an alkyl sulfoacetate surfactant. U.S. Pat. No. 4,113,011 indicates that in a CO.sub.2 foam drive, the problems of low salt tolerance with are typical of both the surface-active sulfates of polyalkoxylated alcohols containing 10 to 16 carbon atoms recommended by U.S. Pat. No. 3,330,346 and the alkyl sulfoacetate surfactants recommended by U.S. Pat. No. 4,088,190 may be avoided by using a surfactant sulfate of a polyalkoxy alcohol containing only 8 or 9 carbon atoms and injecting that surfactant ahead of the CO.sub.2."}
-{"text": "In a telecommunications system where a central telecommunications station supports a plurality of subscribers and a controller is provided for controlling one or more such central telecommunications stations, it is necessary to pass control and other messages between the controller and the central telecommunications station. The messages should be handled in a reliable and efficient manner. It should also be possible to detect messages which are lost and to re-send those messages."}
-{"text": "In recent years, as disclosed in JP 5882522, for example, golf club heads have been proposed in which a raised portion is provided on the crown portion and a sloped surface is formed as a step between the raised portion and the portion rearward thereof. This configuration enables the height of the face portion to be raised by the height of the raised portion. Thus, the rebound performance of the face portion can be improved. Also, on the crown portion, only the raised portion is formed higher, and the portion rearward thereof is formed at a lower position than the raised portion, enabling the center of gravity of the head to be lowered.\nJP 5882522 is an example of related art."}
-{"text": "1. Field of Invention\nThe embodiment of the present invention relates generally to a transmitting method and, more particularly, to a method for transmitting data stream.\n2. Description of Related Art\nIn traditional package transmitting mechanism, when a transmitting end receives an ACK package transmitted from a receiving end, the transmitting end continuously transmits next package. That is to say, the transmitting end stops transmitting packages when the transmitting end does not receive the ACK package transmitted from the receiving end, or the transmitting end disconnects a communication with the receiving end directly when the transmitting end does not receive the ACK package transmitted from the receiving end for a period.\nIn addition, the bandwidth and the buffering space for transmitting the package are different owing to differences of a quality of a content of a film, the way to compress the film, and so on. When a user selects one of films, a client end download related streams of the film from a server. However, in this mode, there will be a bandwidth utilizing shake phenomenon happened in the server and the client end, and the quality of the film will be affected if the package disappears.\nMany efforts have been devoted trying to find a solution of the aforementioned problems. Nonetheless, there still a need to improve the existing apparatus and techniques in the art."}
-{"text": "1. Field of the Invention\nThis invention relates to solar energy assemblies which are adapted to provide decorative and functional means for collecting radiant solar energy and, more specifically, it is directed toward unique panel constructions adapted for such purposes.\n2. Description of the Prior Art\nVarious forms of functional and decorative building construction components positioned on building exteriors such as vertical exterior walls and roofs have been known for years. Not only has it been known to provide decorative wall coverings for the interior, but various forms of exterior siding have been known. See, for example, U.S. Pat. Nos. 2,642,968, 2,777,549, 3,054,223 and 3,394,520.\nAs a result of the shortage of energy on a worldwide basis, more and more effort is being directed toward more efficient use of existing energy supplies. For example, in order to conserve our coal, gas and oil reserves more emphasis has been placed upon maintaining of residential and commercial structures at reduced temperatures in cold weather and providing increased thermal insulation to minimize heat loss. There has also been a great deal of emphasis directed toward the use of solar energy in heating of buildings, heating of hot water and other uses.\nU.S. Pat. No. 3,918,430 discloses a hot water solar system adapted for use on a roof or other portion of a building. A plurality of water channels are housed within a rigid frame underlying a series of layers of plastic material.\nU.S. Pat. No. 4,029,080 discloses a thermal collector for a solar energy system. The prime thrust of this disclosure is directed toward an air system adapted for use on a roof.\nU.S. Pat. No. 4,069,809 discloses a solar system wherein a series of building blocks have transparent members for permitting passage of the sun's rays therethrough. The series of blocks provides a vertical air channel passing immediately behind the transparent window in each block and a series of three generally vertically oriented passageways positioned within each block remote from the front transparent window.\nU.S. Pat. No. 4,120,282 discloses a solar system consisting of a number of fixed flat plate solar reflectors and collectors.\nU.S. Pat. No. 4,073,282 discloses a solar collecting system wherein a matrix of expanded sheets having large openings is employed to collect the sun's radiant energy. Means are provided for circulating air through the chamber and into contact with the slit and expanded sheets.\nU.S. Pat. Nos. 4,076,015 and 4,077,393 each disclose systems wherein modular elements provide a plurality of raised surfaces for receipt of the sun's rays as used in combination with raised reflective surfaces. Among the problems encountered with known solar collecting systems are the somewhat unsightly nature of the same and, in some instances, the expense of installing the same.\nThere remains a need for a solar collecting system for exterior walls, roofs and other portions of buildings which is both decorative and functional. There is a further need for such systems which can be applied readily to preexisting buildings as well as buildings designed and constructed with the solar energy system in mind."}
-{"text": "Prions are infectious pathogens that cause central nervous system spongiform encephalopathies in humans and animals. Prions are distinct from bacteria, viruses and viroids. The predominant hypothesis at present is that no nucleic acid component is necessary for infectivity of prion protein. Further, a prion which infects one species of animal (e.g., a human) will not infect another (e.g., a mouse).\nA major step in the study of prions and the diseases that they cause was the discovery and purification of a protein designated prion protein (\"PrP\") [Bolton et al., Science 218:1309-11 (1982); Prusiner et al., Biochemistry 21:6942-50 (1982); McKinley et al., Cell 35:57-62 (1983)]. Complete prion protein-encoding genes have since been cloned, sequenced and expressed in transgenic animals. PrP.sup.C is encoded by a single-copy host gene [Basler et al., Cell 46:417-28 (1986)] and is normally found at the outer surface of neurons. A leading hypothesis is that prion diseases result from conversion of PrP.sup.C into a modified form called PrP.sup.Sc.\nIt appears that the scrapie isoform of the prion protein (PrP.sup.Sc) is necessary for both the transmission and pathogenesis of the transmissible neurodegenerative diseases of animals and humans. See Prusiner, S. B., \"Molecular biology of prion disease,\" Science 252:1515-1522 (1991). The most common prion diseases of animals are scrapie of sheep and goats and bovine spongiform encephalopathy (BSE) of cattle [Wilesmith, J. and Wells, Microbiol. Immunol. 172:21-38 (1991)]. Four prion diseases of humans have been identified: (1) kuru, (2) Creutzfeldt-Jakob Disease (CJD), (3) Gerstmann-Strassler-Scheinker Disease (GSS), and (4) fatal familial insomnia (FFI) [Gajdusek, D.C., Science 197:943-960 (1977); Medori et al., N. Engl. J. Med. 326:444-449 (1992)]. The presentation of human prion diseases as sporadic, genetic and infectious illnesses initially posed a conundrum which has been explained by the cellular genetic origin of PrP.\nMost CJD cases are sporadic, but about 10-15% are inherited as autosomal dominant disorders that are caused by mutations in the human PrP gene [Hsiao et al., Neurology 40:1820-1827 (1990); Goldfarb et al., Science 258:806-808 (1992); Kitamoto et al., Proc. R. Soc. Lond. 343:391-398. Iatrogenic CJD has been caused by human growth hormone derived from cadaveric pituitaries as well as dura mater grafts [Brown et al., Lancet 340:24-27 (1992)]. Despite numerous attempts to link CJD to an infectious source such as the consumption of scrapie infected sheep meat, none has been identified to date [Harries-Jones et al., J. Neurol. Neurosurg. Psychiatry 51:1113-1119 (1988)] except in cases of iatrogenically induced disease. On the other hand, kuru, which for many decades devastated the Fore and neighboring tribes of the New Guinea highlands, is believed to have been spread by infection during ritualistic cannibalism [Alpers, M. P., Slow Transmissible Diseases of the Nervous System, Vol. 1, S. B. Prusiner and W. J. Hadlow, eds. (New York: Academic Press), pp. 66-90 (1979)].\nThe initial transmission of CJD to experimental primates has a rich history beginning with William Hadlow's recognition of the similarity between kuru and scrapie. In 1959, Hadlow suggested that extracts prepared from patients dying of kuru be inoculated into nonhuman primates and that the animals be observed for disease that was predicted to occur after a prolonged incubation period [Hadlow, W. J., Lancet 2:289-290 (1959)]. Seven years later, Gajdusek, Gibbs and Alpers demonstrated the transmissibility of kuru to chimpanzees after incubation periods ranging form 18 to 21 months [Gajdusek et al., Nature 209:794-796 (1966)]. The similarity of the neuropathology of kuru with that of CJD [Klatzo et al., Lab Invest. 8:799-847 (1959)] prompted similar experiments with chimpanzees and transmissions of disease were reported in 1968 [Gibbs, Jr. et al., Science 161:388-389 (1968)]. Over the last 25 years, about 300 cases of CJD, kuru and GSS have been transmitted to a variety of apes and monkeys.\nThe expense, scarcity and often perceived inhumanity of such experiments have restricted this work and thus limited the accumulation of knowledge. While the most reliable transmission data has been said to emanate from studies using nonhuman primates, some cases of human prion disease have been transmitted to rodents but apparently with less regularity [Gibbs, Jr. et al., Slow Transmissible Diseases of the Nervous System, Vol. 2, S. B. Prusiner and W. J. Hadlow, eds. (New York: Academic Press), pp. 87-110 (1979); Tateishi et al., Prion Diseases of Humans and Animals, Prusiner et al., eds. (London: Ellis Horwood), pp. 129-134 (1992)].\nThe infrequent transmission of human prion disease to rodents has been cited as an example of the \"species barrier\" first described by Pattison in his studies of passaging the scrapie agent between sheep and rodents [Pattison, I. H., NINDB Monograph 2, D. C. Gajdusek, C. J. Gibbs Jr. and M. P. Alpers, eds. (Washington, D.C.: U.S. Government Printing), pp. 249-257 (1965)]. In those investigations, the initial passage of prions from one species to another was associated with a prolonged incubation time with only a few animals developing illness. Subsequent passage in the same species was characterized by all the animals becoming ill after greatly shortened incubation times.\nThe molecular basis for the species barrier between Syrian hamster (SHa) and mouse was shown to reside in the sequence of the PrP gene using transgenic (Tg) mice [Scott et al., Cell 59:847-857 (1989)]. SHaPrP differs from MoPrP at 16 positions out of 254 amino acid residues [Basler et al., Cell 46:417-428 (1986); Locht et al., Proc. Natl. Acad. Sci. USA 83:6372-6376 (1986)]. Tg(SHaPrP) mice expressing SHaPrP had abbreviated incubation times when inoculated with SHa prions. When similar studies were performed with mice expressing the human, or ovine PrP transgenes, the species barrier was not abrogated, i.e., the percentage of animals which became infected were unacceptably low and the incubation times were unacceptably long. Thus, it has not been possible, for example in the case of human prions, to use transgenic animals (such as mice containing a PrP gene of another species) to reliably test a sample to determine if that sample is infected with prions. The seriousness of the health risk resulting from the lack of such a test is exemplified below.\nMore than 45 young adults previously treated with HGH derived from human pituitaries have developed CJD [Koch et al., N. Engl. J. Med. 313:731-733 (1985); Brown et al., Lancet 340:24-27 (1992); Fradkin et al., JAMA 265:880-884 (1991); Buchanan et al., Br. Med. J. 302:824-828 (1991)]. Fortunately, recombinant HGH is now used, although the seemingly remote possibility has been raised that increased expression of wtPrP.sup.C stimulated by high HGH might induce prion disease [Lasmezas et al., Biochem. Biophys. Res. Commun. 196:1163-1169 (1993)]. That the HGH prepared from pituitaries was contaminated with prions is supported by the transmission of prion disease to a monkey 66 months after inoculation with a suspect lot of HGH [Gibbs, Jr. et al., N. Engl. J. Med. 328:358-359 (1993)]. The long incubation times associated with prion diseases will not reveal the full extent of iatrogenic CJD for decades in thousands of people treated with HGH worldwide. Iatrogenic CJD also appears to have developed in four infertile women treated with contaminated human pituitary-derived gonadotrophin hormone [Healy et al., Br. J. Med. 307:517-518 (1993); Cochius et al., Aust. N.Z. J. Med. 20:592-593 (1990); Cochius et al., J. Neurol. Neurosurg. Psychiatry 55:1094-1095 (1992)] as well as at least 11 patients receiving dura mater grafts [Nisbet et al., J. Am. Med. Assoc. 261:1118 (1989); Thadani et al., J. Neurosurg. 69:766-769 (1988); Willison et al., J. Neurosurg. Psychiatric 54:940 (1991); Brown et al., Lancet 340:24-27 (1992)]. These cases of iatrogenic CJD underscore the need for screening pharmaceuticals that might possibly be contaminated with prions.\nRecently, two doctors in France were charged with involuntary manslaughter of a child who had been treated with growth hormones extracted from corpses. The child developed Creutzfeldt-Jakob Disease. (See New Scientist, Jul. 31, 1993, page 4). According to the Pasteur Institute, since 1989 there have been 24 reported cases of CJD in young people who were treated with human growth hormone between 1983 and mid-1985. Fifteen of these children have died. It now appears as though hundreds of children in France have been treated with growth hormone extracted from dead bodies at the risk of developing CJD (see New Scientist, Nov. 20, 1993, page 10.) In view of such, there clearly is a need for a convenient, cost-effective means for removing prions which cause CJD from blood and blood products. The present invention provides such a method."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to securing information in computing systems, and more specifically to limiting access to that information based on the context in which at least a portion of the transactional information was generated, such as from a sale.\n2. Background\nSecuring information has become a priority for organizations to ensure that business processes and information relating thereto remain confidential. As an organization's business processes becomes more complex, the means for securing information has to be flexible to adapt to organizational changes while preserving an appropriate balance between confidentiality (i.e., limiting access to information) and openness (i.e., freedom to access information), both of which are necessary for the success of the organization. Examples of organizational changes requiring such flexibility include employee/group transfers, company reorganizations, compensation plan adjustments and the hiring and/or terminating of personnel.\nTo manage compensation schemes through these types of organizational changes, as well as providing incentive-based compensation for employees in general, organizations have structured compensation plans in accordance with Enterprise Incentive Management (EIM) principles. These principles tailor compensation plans so as to improve optimal performance and to align the organization's strategy with the desired behaviors of it employees. EIM refers generally to managing variable pay plans throughout an organization (i.e., corporation or enterprise) and includes plans for salespeople, suppliers, distribution channel partners, brokers, customers, employees, executives, and partners.\nBut conventional approaches to securing information generated in the framework of an organization typically lack the flexibility to adapt to changes in corporate processes or structure, such as a change in traditional compensation schemes or personnel. For example, consider a personnel change from one part of an organization to another part as shown in FIG. 1.\nFIG. 1 depicts a traditional organizational chart illustrating an employee transferring from one position in organizational structure 100 to another position in new organizational structure 110. Organizational structures 100 and 110 each represent a hierarchical structure depicting supervisor-subordinate relationships where permissions to access secured information decreases from the top position occupied by \u201cA\u201d to the bottom positions occupied by \u201cD,\u201d \u201cE,\u201d and \u201cF.\u201d A square box in FIG. 1, such as the one labeled A, represents a position or role occupied by an employee (or a group/organizational element) and is interrelated with other square boxes, where the interrelations are depicted as lines connecting at least two square boxes. Hence, an employee or organizational element occupying box A is in a supervisory role to employees or organizational elements in boxes \u201cB\u201d and \u201cCynthia,\u201d which are both in subordinate roles to that of box A.\nIn organizational structure 100, Cynthia is shown to occupy a supervisory role in relation to boxes E and F, which can be employees, groups of employees or other organizational elements. In this role, Cynthia has a \u201cspan of access\u201d 102 and is granted permission to access information relating at least to her subordinates occupying boxes E and F, which may include transactional information forming Cynthia's compensation.\nFurther, consider that an employee associated with box E is a sales person operating according to a compensation plan that specifies each of the following allocations to their supervisor Cynthia's compensation: 2% of sales revenue within a particular geographic region; 1% of sales from a particular product line; 2% of sales to a particular customer; and 0.5% of sales by other members of her sales team. Since Cynthia's compensation is based upon such a compensation scheme, she and others in her position are traditionally authorized to access transactional information in her span of access 102 to verify that the sales person's sales revenues are accurately recorded. In particular, Cynthia can access transactional information for boxes E and F but not boxes B and D, which are outside of her span of access 102.\nNext, consider that Cynthia assumes a new role in new organizational structure 110 in the position formerly demarcated as box B of organization structure 100. In this role, traditional security mechanisms allow Cynthia to now have access to transactional information within span of access 112, which authorizes her to examine the activities of the employee(s) relating to box D to review the transactional information that affects her compensation in this new role. But once Cynthia assumes this new role in organizational structure 110, she traditionally is precluded from having access to transactional information for span of access 102. This is generally due to traditional approaches to securing information where a person's set of permissions are dependent on the person's current position in an organization. Without having span of access 102, Cynthia is unable to determine whether her previous efforts and those of her previous subordinates are adequately and accurately recognized so that she can justly be compensated for any activity occurring before she assumed her new role in organization structure 110. Thus, there is a need to provide a flexible method of securing information such that the aforementioned drawbacks of conventional EIM schemes are overcome."}
-{"text": "Montelukast sodium is the active pharmaceutical ingredient of SINGULAIR\u00ae, and is approved for the treatment of asthma and allergic rhinitis. The molecular structure of montelukast is as shown below:\n\nMontelukast sodium is described in U.S. Pat. No. 5,565,473. A crystalline form of montelukast sodium (hereinafter referred to as \u201cForm A\u201d) is described in U.S. Pat. No. 5,614,632."}
-{"text": "1. Field of the Invention\nThis invention relates to a decorative lamp, and particularly to a figurative structure for clamping a decorative lamp string.\n2. Description of the Prior Art\nIn the conventional decorative lamp string for Christmas season, a plurality of sockets are mounted therein by using two or more than two power wires twisted together to connect such sockets in series; such a lamp string is subject to swinging or hanging in the air because of the sockets thereof not being fastened in place.\nIn the conventional decorative lamp string for a given festival, the figurative lamp-mounting frame is usually made of a metal, on which a plurality of socket assemblies connected in series with twisted power wires are mounted thereon. The power wires and the figurative lamp-mounting frame are usually not fastened together; as a result, the lamp string and the lamp-mounting frame are subject to separating from each other. Some of such lamp strings may be fastened in place with fastening cord; however, the sockets and the power wires are also subject to swinging and hanging in the air.\nIn the conventional decorative lamp string for Christmas season, please reference the U.S. Pat. No. 4,802,072, it is a direction fixture for decorative lamp series comprising a socket for a bulb, wires connected to said socket and a retaining ring attached to said socket, said retaining ring being provided with a notch extending longitudinally of the socket with the wires that are connected to said socket being positioned and retained in said socket so as to fixed said socket in a desired orientation and wherein said retaining ring has an outer face which is in registry with an end rim of said socket. In the aforesaid invention, the socket assemblies and the power wires are fastened in place with fastening slips, but the socket assemblies and the power wires appear to be out of order; therefore, the lamp string has to be fastened to a lamp-mounting frame with fastening cords.\nIn another conventional lamp string for Christmas season, please reference the U.S. Pat. No. 5,526,246, the front part of the figurative structure is provided with a fastening base, which is furnished with a plurality of hooks for clamping sockets respectively; then, the sockets are fastened to the figurative structure. When the bulbs wink, the figurative structure will be shown vividly. The back of the figurative structure is provided with a connection plate; as a result, such figurative structure can only be used for one-side decoration.\nIn still another conventional lamp string for Christmas season, such as U.S. Pat. No. 5,727,872, the bulb-plugging end of the socket assembly has two curved surfaces on both sides thereof, and one side thereof has two symmetrical arm plates extended out; the tail ends of such arm plates have two curved hooks respectively; the inner surface of the hooks each have a curved surface to fit to the figurative lamp-mounting frame. The outer ends of the two hooks form into an opening; each socket assembly has two arm plates to form into an opening so as to click to the metal rod of the figurative lamp-mounting frame, i.e., to have the socket assembly clamped to the metal rod of the figurative lamp-mounting frame."}
-{"text": "1. Field of the Invention\nThe present invention relates to a pneumatic gripper comprising at least one pneumatic structural element.\n2. History of the Related Art\nThe devices nearest to the present invention are known from U.S. Pat. No. 3,056,625, Timmerman (D1) and JP 05261687, Bridgestone (D2).\nD1 describes a gripper for goods which is configured as a clamp and has grippers held movably on its lateral vertically disposed sections by means of hinges. Located between each vertically disposed section and the associated gripper is an inflatable bellows which, under pressure, pushes the gripper away from the vertical section towards the inside of the gripper so that the goods are grasped by the gripper.\nD2 also describes a gripper configured in the manner of a clamp with lateral vertically disposed sections. A supporting arm runs parallel to each of the vertical sections, on the inner side thereof, said supporting arm for its part being movably hinged in the cross member of the gripper by means of a hinge. Inflatable cushions are provided on each supporting arm which, when filled with compressed air, press the supporting arms away towards the inside (whereby said cushions also move inwards since they are disposed on the supporting arms).\nThese known grippers have the disadvantage of a rigid structure, for example, with the consequence that they can only be used in a vertical position."}
-{"text": "1. Field of the Invention\nThe present invention relates to the game of pool.\nMore particularly, the present invention relates to a portable pool game.\n3. Description of the Prior Art\nEvery pool player, whether because of playing a better player or because of playing someone on a run, has unhappily experienced watching the game more than shooting in the game. In a conventional game of pool, players shot until they missed. If a player was good he could end up shooting for a long period of time.\nNow through the games of \"Super Pool\", and \"Hourglass Pool\", of the present invention, all of the old games, and the new games introduced through the addition of six new balls, can be played using equal shooting time, if desired. The option to play equal shooting time games is still present, if desired.\nNumerous innovations for the game of pool have been provided in the prior art that are adapted to be used. Even though these innovations may be suitable for the specific individual purposes to which they address, they would not be suitable for the purposes of the present invention as heretofore described."}
-{"text": "Coronary artery disease may produce coronary lesions in the blood vessels providing blood to the heart, such as a stenosis (abnormal narrowing of a blood vessel). As a result, blood flow to the heart may be restricted. A patient suffering from coronary artery disease may experience chest pain, referred to as chronic stable angina during physical exertion or unstable angina when the patient is at rest. A more severe manifestation of disease may lead to myocardial infarction, or heart attack.\nA need exists to provide more accurate data relating to coronary lesions, e.g., size, shape, location, functional significance (e.g., whether the lesion impacts blood flow), etc. Patients suffering from chest pain and/or exhibiting symptoms of coronary artery disease may be subjected to one or more tests that may provide some indirect evidence relating to coronary lesions. For example, noninvasive tests may include electrocardiograms, biomarker evaluation from blood tests, treadmill tests, echocardiography, single positron emission computed tomography (SPECT), and positron emission tomography (PET). These noninvasive tests, however, typically do not provide a direct assessment of coronary lesions or assess blood flow rates. The noninvasive tests may provide indirect evidence of coronary lesions by looking for changes in electrical activity of the heart (e.g., using electrocardiography (ECG)), motion of the myocardium (e.g., using stress echocardiography), perfusion of the myocardium (e.g., using PET or SPECT), or metabolic changes (e.g., using biomarkers).\nFor example, anatomic data may be obtained noninvasively using coronary computed tomographic angiography (CCTA). CCTA may be used for imaging of patients with chest pain and involves using computed tomography (CT) technology to image the heart and the coronary arteries following an intravenous infusion of a contrast agent. However, CCTA also cannot provide direct information on the functional significance of coronary lesions, e.g., whether the lesions affect blood flow. In addition, since CCTA is purely a diagnostic test, it cannot be used to predict changes in coronary blood flow, pressure, or myocardial perfusion under other physiologic states, e.g., exercise, nor can it be used to predict outcomes of interventions.\nThus, patients may also require an invasive test, such as diagnostic cardiac catheterization, to visualize coronary lesions. Diagnostic cardiac catheterization may include performing conventional coronary angiography (CCA) to gather anatomic data on coronary lesions by providing a doctor with an image of the size and shape of the arteries. CCA, however, does not provide data for assessing the functional significance of coronary lesions. For example, a doctor may not be able to diagnose whether a coronary lesion is harmful without determining whether the lesion is functionally significant. Thus, CCA has led to what has been referred to as an \u201coculostenotic reflex\u201d of some interventional cardiologists to insert a stent for every lesion found with CCA regardless of whether the lesion is functionally significant. As a result, CCA may lead to unnecessary operations on the patient, which may pose added risks to patients and may result in unnecessary heath care costs for patients.\nDuring diagnostic cardiac catheterization, the functional significance of a coronary lesion may be assessed invasively by measuring the fractional flow reserve (FFR) of an observed lesion. FFR is defined as the ratio of the mean blood pressure downstream of a lesion divided by the mean blood pressure upstream from the lesion, e.g., the aortic pressure, under conditions of increased coronary blood flow, e.g., induced by intravenous administration of adenosine. The blood pressures may be measured by inserting a pressure wire into the patient. Thus, the decision to treat a lesion based on the determined FFR may be made after the initial cost and risk of diagnostic cardiac catheterization has already been incurred.\nThus, a need exists for a method for assessing coronary anatomy, myocardial perfusion, and coronary artery flow noninvasively. Such a method and system may benefit cardiologists who diagnose and plan treatments for patients with suspected coronary artery disease. In addition, a need exists for a method to predict coronary artery flow and myocardial perfusion under conditions that cannot be directly measured, e.g., exercise, and to predict outcomes of medical, interventional, and surgical treatments on coronary artery blood flow and myocardial perfusion.\nIt is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory only and are not restrictive of the disclosure."}
-{"text": "1. Field of the Invention\nThe present invention relates to an optical system of an optical pickup in an optical information recording and reproducing apparatus which records and reproduces information for optical discs having different corresponding wavelengths. More particularly, the present invention relates to an optical information recording and reproducing apparatus which allows compatibility for a plurality of optical recording mediums using laser light sources of different wavelengths, to an optical pickup, to an objective lens module, and to a diffractive optical element.\n2. Description of Related Art\nAs an optical information recording and reproducing apparatus, an optical disc apparatus is known in which recorded information can be read from an optical recording medium, that is, an optical disc, such as digital versatile disc (hereinafter, referred to as DVD), compact disc (hereinafter, referred to as CD), or the like.\nA compatible optical disc apparatus is known in which recorded information can be read from DVD and CD. As for DVD, the substrate thickness is 0.6 mm, the corresponding wavelength is in a range of 635 nm to 655 nm, and the numerical aperture (NA) of an objective lens is about 0.6. As for CD, the substrate thickness is 1.2 mm, the corresponding wavelength is in a range of 760 to 800 nm, and the numerical aperture of an objective lens is about 0.45. In the compatible optical disc apparatus, there is a case in which a laser light source having a wavelength \u03bbDVD in the vicinity of the wavelength 660 nm for DVD and a laser light source having a wavelength \u03bbCD in the vicinity of the wavelength 780 nm for CD are mounted.\nFor example, a technology is suggested in which an optical pickup device for allowing information to be recorded and reproduced for information recording mediums having different substrate thicknesses for DVD/CD, and an objective lens and an optical element used for the optical pickup device are provided (JP-A-2001-235676). The optical pickup device is suggested in which the objective lens having diffractive orbicular zones is used for the optical pickup device, such that, with an outside light flux of a predetermined numerical aperture in a use state of a small numerical aperture as a flare, recording and reproducing of information are performed for various information recording mediums having different thicknesses. The objective lens having such diffractive orbicular zones includes a diffraction surface having the diffractive orbicular zones. Here, when a function of an optical path difference of the diffraction surface is \u03a6(h) (where h is a distance from an optical axis), d\u03a6(h)/dh is a discontinuous or substantially discontinuous function at a place of a predetermined distance h.\nOn the other hand, as for blue-ray disc (hereinafter, referred to as BD), the thickness of a transmissive protection layer (which corresponds to the thickness of a transparent substrate of DVD or the like) is 0.1 mm, the corresponding wavelength is 408 nm, and the numerical aperture of an objective lens is about 0.85. Accordingly, in a BD/DVD/CD compatible optical disc apparatus, a laser light source which emits laser light of \u03bbBD in the vicinity of the wavelength 408 nm, that is, an optical system, needs to be mounted, in addition to the configuration of the above-described compatible optical disc apparatus. Further, since the optical discs of BD, DVD, and CD have different thicknesses, a unit for correcting three kinds of different spherical aberrations needs to be provided. In addition, since all of them have different numerical apertures, a corresponding unit also needs to be provided. However, in JP-A-2001-235676 described above, the specified descriptions of these units are not given. That is, it is difficult to realize compatibility of three or more kinds of recording mediums having different light source wavelengths, numerical apertures (effective diameters), optical disc thicknesses (the thickness of a transmissive protection layer), such as BD, DVD, CD, and the like by use of a single objective lens according to the related art.\nIn order to realize an optical pickup for a compatible apparatus, a method is suggested in which an objective lens exclusively used for BD and a DVD/CD compatible objective lens are used, and are switched according to wavelengths. In this case, however, since two objective lenses are used, a complex lens switching mechanism needs to be provided, which causes a problem in that manufacturing costs are increased. In addition, since an actuator is made large, it is disadvantageous to reduce the size of the apparatus. Further, a method may be considered in which an objective lens and a collimator lens are incorporated, but, since the collimator is fixed with respect to the objective lens, it may be difficult to maintain performance at the time of movement of the objective lens.\nIn any cases, if a plurality of light sources are used, and an optical system of exclusive prism, lens, and the like is configured in order to ensure compatibility of BD, DVD, and CD, an optical pickup or an overall optical head is complicated, and tends to have the large size."}
-{"text": "The Internet has become a popular place to conduct business. Through Web auction sites, Web sites for displaying classified ads, Web shopping malls, online chat rooms, and other online transaction facilitation sites, two consumers may agree to a transaction. Frequently, such transactions involve the exchange of goods or services for money. While consumers frequently find that agreeing to transactions on the Internet is easy, completing a payment to consummate the transaction is more difficult.\nTypically, two consumers who have agreed through the Internet to exchange goods for money resort to offline methods to perform the exchange. For example, the seller may ship the goods to the buyer through a shipping service, and the buyer may send a paper check to the seller.\nSuch offline methods of exchange are problematic. Because the buyer and the seller are usually strangers, they may not trust each other to perform their mutual obligations under the agreement. Accordingly, they may be unable to agree whether the buyer will send the check first or the seller will send the goods first. Even if the buyer and the seller agree that the seller will ship the goods at the same time as the buyer sends the check, the seller has no guarantee that the check will not bounce. Likewise, the buyer has no guarantee that the goods will arrive in satisfactory condition. Accordingly, a significant percentage of transactions to which an individual buyer and seller have agreed upon over the Internet are never consummated.\nAnother inconvenience of transactions agreed upon by individuals over the Internet is that the buyer is often limited to paying by cash or paper check. More convenient payment instruments exist, such as credit cards and bank account debits through electronic fund transactions. However, the buyer typically does not have the option to use these other payment instruments when the seller is an individual as opposed to a retail business that has been pre-established as an online merchant.\nThe term \u201cmerchant\u201d is used herein to refer to a seller of goods or services who is authorized by a credit card association (such as DISCOVER, VISA, or MASTERCARD) to submit to the credit card association charges on credit cards belonging to members of the credit card association. After receiving an authorization for the charge, the merchant then receives from the credit card association a direct deposit into the merchant's bank account of the amount of the charge. As known to those skilled in the art, a business must undergo an approval process in order to become a merchant, and upon approval, the merchant is assigned a merchant number.\nAlthough retail businesses are routinely set up as merchants in order to accept payments through credit cards or electronic fund transactions, this is not an adequate solution to facilitating payments between individuals over the Internet. For example, merchants, after undergoing an extensive underwriting effort, are typically given special privileges, such as a general authorization to charge credit cards. This general authorization provides the merchant with the ability to commit fraud. Specifically, the merchant is capable of charging a customer's credit card more than he should. Also, the merchant may submit charges on a credit card belonging to a credit card association member with which the merchant has never had any contact. For these reasons, the idea of allowing individual sellers to become merchants has heretofore been rejected.\nAnother problem with an individual seller becoming a merchant is that the approval process for becoming a merchant is frequently more of a hassle than an occasional seller is willing to undergo. The purpose of the approval process is to reduce the risk of fraud by the merchant. Accordingly, the seller usually must submit extensive background information for consideration in the approval process. This may be inconvenient and time consuming for the seller.\nTherefore, there is a need in the art for a safe and convenient method by which one consumer can pay a second consumer over the Internet."}
-{"text": "The Kiosk was designed and developed to accommodate a need for a \u201cstand alone unit\u201d that houses an interactive computer monitor/touch screen display, for commercial trade shows and traveling exhibit applications. The requirements were to be lightweight, collapsible and shippable (UPS, FEDEX etc.) and yet maintain a \u201ccorporate\u201d look. It was also required to have shelving for a CPU and CD/DVD player, keyboard, speakers for the audio aspect and a storage area for miscellaneous accessories. The kiosk also required that an operator be able to gain access to the equipment without completely disassembling the unit, so a locking door feature was incorporated into the present invention.\nFrom the foregoing, it will be appreciated that there is a need in the art to develop a sturdy lightweight foldable monitor stand with foldable shelves for easy portability and storage. The present invention is directed to overcoming one, or more, of the problem set forth above."}
-{"text": "1. Field of the Invention\nThe present invention relates to a technique for collectively connecting a plurality of grounding wires included in a wire harness for a vehicle to a given ground site inside the vehicle.\n2. Description of the Related Art\nHeretofore, as a ground connecting device for collectively connecting a plurality of grounding wires included in a wire harness for a vehicle, to a ground site of the vehicle, there has been known a type described in JP 10-208815A.\nFIG. 10 shows an outline of this device. The device comprises a harness-side connector 7 to be provided at a terminal end of a wire harness including a plurality of grounding wires, and a ground joint connector 1 to be fixed to a given ground site (in FIG. 10, a bolt 6) provided on a vehicle body 3. The harness-side connector 7 includes a plurality of non-illustrated female terminals to be attached to respective terminal ends of the grounding wires and a connector housing 8 for collectively holding the female terminals. The harness-side connector housing 8 has a plurality of built-in terminal locking portions for holding the female terminals respectively. The ground joint connector 1 includes a grounding conductor 5 and a connector housing 2 which holds the grounding conductor 5, the grounding conductor 5 integrally having a grounding terminal portion 4 to be fixed to the ground site and a plurality of non-illustrated male terminals provided inside the connector housing 2.\nAccording to this device, interconnecting the ground joint connector 1 and the harness-side connector 7 and fixing the grounding terminal portion 4 in the ground joint connector to the bolt 6 as the ground site establish a collective connection of the grounding wires to the ground site. Specifically, the female terminals held by the connector housing 8 of the harness-side connector 7 and the male terminals of the grounding conductor 5 held by the connector housing 2 of the ground joint connector 1 are fitted to each other respectively, thus electrically connecting the grounding wires to which the female terminals are attached to the ground site through the female terminals and the grounding conductor 5; simultaneously, the connector housing 8 of the harness-side connector 7 and the connector housing 2 of the ground joint connector 1 are fitted to each other, and this fitting is locked by engagement between respective engagement portions provided in the two connector housings 8 and 2, the lock keeping the female and male terminals fitted to each other respectively.\nHowever, this ground connecting device, occupying a large space, is difficult to use in a little space in a vehicle. Specifically, the harness-side connector 7 and the ground joint connector 1 of the device require the connector housings 8 and 2 for holding the terminals respectively; furthermore, the connector housings 8 and 2 occupy a large space as a whole for their mutual fitting and the lock of the fitting. To avoid interference between the connector housings 8 and 2 and the vehicle body 3, the connectors 7 and 1 are required to protrude in a large size from an inner surface of the vehicle body 3. Particularly, the case of connecting a grounding terminal 9 attached to an extra grounding wire W to the grounding terminal portion 4 so as to superimpose them to each other as shown in FIG. 10 requires a large gap size L between the vehicle body 3 and each of the connector housings 8 and 2 as shown in FIG. 10, in order to avoid the interference between the grounding terminal and each of the connector housings 8 and 2; this causes the entire device to occupy a larger space."}
-{"text": "1. Field of the Invention\nThe present invention pertains to typing recognition systems and methods, and more particularly to recognition of typing in air or on a relatively smooth surface that provides less tactile feedback than conventional mechanical keyboards.\n2. The Related Art\nTypists generally employ various combinations of two typing techniques: hunt and peck and touch typing. When hunting and pecking, the typist visually searches for the key center and strikes the key with the index or middle finger. When touch typing, the fingers initially rest on home row keys, each finger is responsible for striking a certain column of keys and the typist is discouraged from looking down at the keys. The contours and depression of mechanical keys provide strong tactile feedback that helps typists keep their fingers aligned with the key layout. The finger motions of touch typists are ballistic rather than guided by a slow visual search, making touch typing faster than hunt and peck. However, even skilled touch typists occasionally fall back on hunt and peck to find rarely-used punctuation or command keys at the periphery of the key layout.\nMany touchscreen devices display pop-up or soft keyboards meant to be activated by lightly tapping a displayed button or key symbol with a finger or stylus. Touch typing is considered impractical on such devices for several reasons: a shrunken key layout may have a key spacing too small for each finger to be aligned with its own key column, the smooth screen surface provides no tactile feedback of finger/key alignment as keys are struck, and most touchscreens cannot accurately report finger positions when touched by more than one finger at a time. Such temporal touch overlap often occurs when typing a quick burst of keys with both hands, holding the finger on modifier keys while striking normal keys, or attempting to rest the hands. Thus users of touchscreen key layouts have had to fall back on a slow, visual search for one key at a time.\nSince touchscreen and touch keyboard users are expected to visually aim for the center of each key, typing recognition software for touch surfaces can use one of two simple, nearly equivalent methods to decide which key is being touched. Like the present invention, these methods apply to devices that report touch coordinates interpolated over a fine grid of sensors rather than devices that place a single large sensor under the center of each key. In the first method, described in U.S. patent application Ser. No. 09/236,513 by Westerman and Elias, the system computes for each key the distance from key center to the sensed touch location. The software then selects the key nearest the finger touch. In the second method, described in U.S. Pat. No. 5,463,388 to Boie et al., the software establishes a rectangle or bounding box around each key and decides which, if any, bounding box the reported touch coordinates lie within. The former method requires less computation, and the latter method allows simpler control over individual key shape and guard bands between keys, but both methods essentially report the key nearest to the finger touch, independent of past touches. Hence we refer to them as \u2018nearest key\u2019 recognizers.\nUnlike touchscreens, the multi-touch surface (MTS) described by Westerman and Elias in Ser. No. 09/236,513 can handle resting hands and temporal finger overlap during quick typing bursts. Since the MTS sensing technology is fully scalable, an MTS can easily be built large enough for a full-size QWERTY key layout. The only remaining barrier to fast touch typing on an MTS is the lack of tactile feedback. While it is possible to add either textures or compressibility to an MTS to enhance tactile feedback, there are two good reasons to keep the surface firm and smooth. First, any textures added to the surface to indicate key centers can potentially interfere with smooth sliding across the surface during multi-finger pointing and dragging operations. Second, the MTS proximity sensors actually allow zero-force typing by sensing the presence of a fingertip on the surface whether or not the finger applies noticeable downward pressure to the surface. Zero-force typing reduces the strain on finger muscles and tendons as each key is touched.\nWithout rich tactile feedback, the hands and individual fingers of an MTS touch typist tend to drift out of perfect alignment with the keys. Typists can limit the hand drift by anchoring their palms in home position on the surface, but many keystrokes will still be slightly off center due to drift and reach errors by individual fingers. Such hand drift and erroneous finger placements wreak havoc with the simple \u2018nearest key\u2019 recognizers disclosed in the related touchscreen and touch keyboard art. For example, if the hand alignment with respect to the key layout drifts by half a key-spacing (\u02dc9 mm or \u215c\u2033), all keystrokes may land half-way between adjacent keys. A \u2018nearest key\u2019 recognizer is left to choose one of the two adjacent keys essentially at random, recognizing only 50% of the keystrokes correctly. A spelling model integrated into the recognizer can help assuming the typist intended to enter a dictionary word, but then actually hinders entry of other strings. Thus there exists a need in the touchscreen and touch keyboard art for typing recognition methods that are less sensitive to the hand drift and finger placement errors that occur without strong tactile feedback from key centers.\nFor many years, speech, handwriting, and optical character recognition systems have employed spelling or language models to help guess users' intended words when speech, handwriting, or other input is ambiguous. For example, in U.S. Pat. No. 5,812,698 Platt et al. teach a handwriting recognizer that analyzes pen strokes to create a list of probable character strings and then invokes a Markov language model and spelling dictionary to pick the most common English word from that list of potential strings. However, such systems have a major weakness. They assume all user input will be a word contained in their spelling or language model, actually impeding entry of words not anticipated by the model. Even if the user intentionally and unambiguously enters a random character string or foreign word not found in the system vocabulary, the system tries to interpret that input as one of its vocabulary words. The typical solution is to provide the user an alternative (often comparatively clumsy) process with which to enter or select strings outside the system vocabulary. For example, U.S. Pat. No. 5,818,437 to Grover et al. teaches use of a dictionary and vocabulary models to disambiguate text entered on a \u2018reduced\u2019 keyboard such as a telephone keypad that assigns multiple characters to each physical key. In cases that the most common dictionary word matching an input key sequence is not the desired word, users must select from a list of alternate strings. Likewise, users of speech recognition system typically fall back on a keyboard to enter words missing from the system's vocabulary.\nUnfortunately, heavy reliance on spelling models and alternative entry processes is simply impractical for a general-purpose typing recognizer. Typing, after all, is the fallback entry process for many handwriting and speech recognition systems, and the only fallback conceivable for typing is a slower, clumsier typing mode. Likewise, personal computer users have to type into a wide variety of applications requiring strange character strings like passwords, filenames, abbreviated commands, and programming variable names. To avoid annoying the user with frequent corrections or dictionary additions, spelling model influence must be weak enough that strings missing from it will always be accepted when typed at moderate speed with reasonable care. Thus a general-purpose typing recognizer should only rely on spelling models as a last resort, when all possible measurements of the actual typing are ambiguous."}
-{"text": "1. Field of the Invention\nThe present invention relates to a power driven screwdriver having a clutch mechanism for transmitting rotation of a drive motor to a spindle with a driver bit.\n2. Description of the Prior Art\nIn a power driven screwdriver, a clutch mechanism is provided for transmitting and disconnecting the rotation of a drive motor to a spindle with a driver bit. The clutch mechanism is normally constructed as a claw clutch and includes a pair of clutch members, one of which is mounted on the spindle and the other of which is mounted on a main gear driven by the drive motor. The spindle is movable in an axial direction for engaging and disengaging the clutch members. With such a clutch mechanism constructed by a simple claw clutch, since the rotation of the spindle is restrained, for example, at the completion of a screw driving operation, the clutch mechanism temporarily repeats its engaging and disengaging operation. This will generate clanging sounds, giving unpleasant feeling to the operator, and cause early wear of the clutch mechanism.\nU.S. Pat. No. 4,655,103 discloses a power driven screwdriver including stopper for adjusting the driving\na amount of a screw by a driver bit. A claw clutch mechanism is provided between a driver shaft and a spindle movable in an axial direction. The claw clutch mechanism includes a first and a second clutch member formed on the driver shaft and the spindle, respectively. A clutch disc is interposed between the driver shaft and the spindle and includes a third and a fourth clutch member for engagement with the first and second clutch members respectively. A spring is interposed between the first and third clutch members for normally keeping them at a disengaging position. The second and fourth clutch members includes relief portions which serves not to transmit rotation. When the stopper abuts on a work to be screwed, the driver shaft continues rotation while the rotation of the spindle is prevented. This may cause the operation of the relief portions of the second and fourth clutch members to positively disengage the first and the third clutch members with the aid of the spring.\nU.S. Pat. No. 4,809,572 discloses a power driven screwdriver including a stopper sleeve for adjusting the driving amount of a screw and a claw clutch mechanism having a pair of clutch members, one of which is mounted on a main gear driven by a drive motor, while the other of which is mounted on a spindle. A spring is provided for normally keeping the clutch member of the spindle out of engagement with the clutch member of the main gear. A control mechanism is provided between the spindle and the clutch member mounted on the spindle. The control mechanism includes oblique recesses and a ball for engagement with the recesses. With such construction, when the stopper sleeve abuts on a work to be screwed, the main gear continues its rotation while the rotation of the spindle is prevented. In this stage, the control mechanism operates to positively move the clutch member of the spindle out of engagement with the clutch member of the main gear with the aid of the spring.\nHowever, with the above prior U.S. Patents, the operation of the clutch mechanism must accompany a reciprocal movement of the spindle at a long distance. In general, a power driven screwdriver is provided with a seal member for sealing between a spindle and a housing to prevent entry of dust within the housing. In case the spindle reciprocally moves at a long distance, the dust may be absorbed into the housing through the ga between the seal member and the spindle or the housing by the pumping effect. Thus, when the spindle moves into the housing, negative pressure will be created in the housing. Such dust entered into the housing may cause early wear or damage of the clutch mechanism or bearings disposed within the housing.\nFurther, with the clutch mechanism of the above U.S. Patents, after the stopper or the stopper sleeve has abutted on the work, no further driving operation cannot be made even if the driving of a screw wa insufficient."}
-{"text": "Such a cartridge is known from DE 10 2008 057 443 A1, where the functional element is a valve device.\nIt is known from U.S. Pat. No. 4,391,590 B1 that a functional element is designed as a cap, which is pulled over a cannula opening of a cannula duct to close the cannula section. It is disadvantageous here that the cartridge and the cap for the cartridge must be manufactured in two mutually independent manufacturing steps, wherein the cap is placed, as a rule, by hand by a person or in an automated manner on the cannula opening of the cannula section. This causes high manufacturing costs. In addition, with the cap already placed, there is a risk of air inclusions during the filling of the cartridge with the dental material, as result of which the shelf life of the dental material may be reduced and/or the quantity of the filling may show undesired variations in a comparison of a plurality of cartridges. Such air inclusions may lead to the loss of the cap, especially during transportation, because of the expansion of the air, as a result of which the storage stability is reduced. Even though the inclusion of air can be reduced when filling the cartridge without cap, there is a risk now that dental material will escape from the cannula opening of the cannula section, as a result of which dental material will be lost. This leads to higher manufacturing costs. Such cartridges are intended for single-time use especially in the field of dentistry. However, there is a risk when using caps that the cannula section will be reclosed with the cap in order to use a residual material that is preset later. There is a risk of contamination of the dental material and/or of an increased risk for infection because of its undesired reclosing of the cartridge that was once opened.\nA cartridge, in which fibers or a flocking are connected to the cartridge in the area of an outlet of the cannula section, is known from U.S. Pat. No. 6,059,570. This functional element is used to apply, spread and/or burnish the dental material. It is disadvantageous here that the application of the fibers or of the flocking is carried out in an independent manufacturing step and fully independently from the manufacture of the cartridge. It is also disadvantageous that the cannula section is rigid in the area of the fibers or flocking. As a result, there is a risk that a treatment with the functional element is perceived by a patient as being unpleasant and/or painful. In addition, there is a risk that undesired injuries will develop because of the rigid design. The spreading and/or burnishing of the dental material is also made difficult by the rigid design of the cartridge and of the cannula section.\nFurthermore, it is disadvantageous in prior-art cartridges that there is a risk of an especially abrupt rupture of material in case of an overstressing due to an excessively strong force or an excessively high pressure being applied to press the dental material out of the cartridge and/or the reservoir."}
-{"text": "This invention relates to an infrared detecting element and also an infrared imaging device.\nSome infrared detectors use Si crystals and detect infrared rays having wavelengths equal to or longer than several micrometers. Such infrared detectors are of two types, the first type being produced by doping impurities into the Si crystals and the second type using heterojunction barriers.\nInfrared Detectors II, Chapter 2, Semiconductors and Semimetals, written by P. R. Bratt, published from Academic Press in 1977, discloses the first-type infrared detectors.\nJapanese published unexamined patent application 61-241985 discloses the second-type infrared detector. The documents \"3P79\" of the lecture in the thirty-third spring meeting of Applied Physical Society of Japan in 1986 also discloses the second-type infrared detector.\nThese two types of infrared detectors are useful for infrared two-dimensional imaging devices of a monolithic type. The first-type infrared detectors have the following drawback. Since the quantity of doped impurities is limited, the detector sensitivity is low and the detected wavelength is fixed in dependence on the type of the impurities. Accordingly, it is impossible to maximize the detector sensitivity at an arbitrary wavelength. The second-type infrared detectors are free from such a drawback."}
-{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus such as a copying machine, a facsimile machine, a printer and a multifunction machine, and, moreover, to a sheet stacking device stacking a sheet (recording medium) formed with an image, and a sheet processing device performing post processing of a sheet.\n2. Description of the Related Art\nA sheet processing device with the following configuration has been well known as a sheet processing device into which a sheet with an image formed in an image forming apparatus is conveyed. The sheet processing device has a buffer roller through which, when the sheet processing device receives sheets, which have been formed with an image, and have been discharged from an image forming apparatus main body, the received sheets are superimposed for temporary waiting before the sheets are conveyed to a post processing mechanism such as a stapling machine and a saddle stitching machine. That is, while a preceding sheet bundle is processed in a processing tray, first several sheets of the succeeding sheet bundle are on the buffer roller for temporary waiting. When the preceding sheet bundle, which has been processed, is discharged from the processing tray, the succeeding several sheets, which have been delivered from the buffer roller, are conveyed to the processing tray. A brief explanation of a sheet post-processing device with the above-described configuration will be given, referring to FIG. 9A through FIG. 9C.\nA plurality of sheets P1, P2, . . . are superimposed one on top of another and wrapped around a buffer roller 5 to form a wrapping path 32. For example, three sheets P1, P2, and P3 are superimposed one on top of another, delivered from the path 32 after temporary waiting, conveyed, and conveyed to a processing tray 101 through a discharge roller 7, bundle discharge rollers 180a, and 180b. When the rear ends of the sheets pass the discharge roller 7, the bundle discharge rollers 180a and 180b rotate in the reverse direction in such a way that the sheet bundle of three sheets P1, P2, and P3 is returned in the direction in which the sheets abut against a rear-end stopper 3 of the processing tray 101. Alignment is performed in such a way that the bundle discharge roller 180b is separated from the bundle discharge roller 180a just before the rear end of the sheet bundle abuts against the rear-end stopper 3 and the sheet bundle abuts against the rear-end stopper 3 by moving inertia. At this time, alignment in a direction perpendicular to the conveyance direction is performed, using aligning plates.\nWhen all the sheets of the first sheet bundle are aligned on the processing tray 101 in such a manner, a swinging guide 150 is lowered and the bundle discharge roller 180b sits atop the sheet bundle to perform stitching processing of the sheet bundle, and the like, using a processing machine such as a stapling machine indicated by a reference number 4 in FIG. 9A.\nAccording to the above-described procedures, a first plurality of sheets of the succeeding second sheet bundle are wrapped around the buffer roller 5 as a temporary accumulating unit for waiting until processing for the first sheet bundle is completed. Thereby, a high-speed image forming apparatus by which sheets are discharged from the main body of an image forming apparatus at a small interval may be realized. Incidentally, the varieties of the quality and the size of sheets have been further increased in recent years. But the sheet processing device shown in FIG. 9A through FIG. 9C may hardly treat sheets, for example, special paper such as coated paper, the surface of which is treated, thick one, and large-sized one.\nEven if these kinds of sheets may be surely aligned one by one, it is difficult to align a plurality of the sheets in a state in which the sheets are superimposed. That is, the rear ends of a plurality of the sheets with a special quality, or with a special sheet size is run into the rear-end stopper 3 on the processing tray 1. In this case, there is generated a state in which all the three sheets P1, P2, and P3 with a large coefficient of friction, such as that of coated paper, are not completely returned to the rear end of the stopper 3. Especially, it is serious that the sheet P2 such as the second sheet of the superimposed ones is incompletely or faultily returned, that is, the quality of the post processing is deteriorated, and, consequently, the productivity is reduced. Moreover, when a plurality of sheets such as a thick one, and a large-sized one are superimposed and moved, there is caused larger inertia than the one caused in a case in which one sheet is moved. Accordingly, there is a case in which non-aligning is caused, because the sheets are vigorously run into the rear-end stopper 3 and bound. Moreover, there is a possibility that the end portion of the sheet buckles, and is damaged."}
-{"text": "This invention relates to a reference clock architecture for an integrated circuit device, and particularly for types of integrated circuit devices, such as programmable devices, where a user may specify a clock rate.\nCertain types of integrated circuit devices allow users to specify various settings, such as clock rates. In particular, programmable devices, including, for example, programmable logic devices such as field-programmable gate arrays (FPGAs), may allow a user to specify a complete logic configuration, various portions of which may require different clock rates, none of which are known with any certainty at the time of device manufacture. Such devices have been manufactured with circuitry to allow various clock rates to be selected by the user, which may have resulted in overly complex clock networks, including many components that may never be used by a particular user.\nFor example, such devices may incorporate high-speed serial interfaces to accommodate high-speed (i.e., greater than 1 Gbps) serial I/O standards. Because there are multiple different standards, which may operate at multiple different rates, and because a user may elect to use more than one standard and/or rate, the ability to provide multiple reference clocks may be desirable. Heretofore, this has required the provision of multiple reference clock sources such as phase-locked loops (PLLs) or delay-locked loops (DLLs), with a clock network capable of routing a reference clock signal from any one of those sources to any one of a number of interface circuits."}
-{"text": "Technical Field of the Invention\nThis invention relates generally to computing systems and more particularly to data storage solutions within such computing systems.\nDescription of Related Art\nComputers are known to communicate, process, and store data. Such computers range from wireless smart phones to data centers that support millions of web searches, stock trades, or on-line purchases every day. In general, a computing system generates data and/or manipulates data from one form into another. For instance, an image sensor of the computing system generates raw picture data and, using an image compression program (e.g., JPEG, MPEG, etc.), the computing system manipulates the raw picture data into a standardized compressed image.\nWith continued advances in processing speed and communication speed, computers are capable of processing real time multimedia data for applications ranging from simple voice communications to streaming high definition video. As such, general-purpose information appliances are replacing purpose-built communications devices (e.g., a telephone). For example, smart phones can support telephony communications but they are also capable of text messaging and accessing the internet to perform functions including email, web browsing, remote applications access, and media communications (e.g., telephony voice, image transfer, music files, video files, real time video streaming. etc.).\nEach type of computer is constructed and operates in accordance with one or more communication, processing, and storage standards. As a result of standardization and with advances in technology, more and more information content is being converted into digital formats. For example, more digital cameras are now being sold than film cameras, thus producing more digital pictures. As another example, web-based programming is becoming an alternative to over the air television broadcasts and/or cable broadcasts. As further examples, papers, books, video entertainment, home video, etc. are now being stored digitally, which increases the demand on the storage function of computers.\nA typical computer storage system includes one or more memory devices aligned with the needs of the various operational aspects of the computer's processing and communication functions. Generally, the immediacy of access dictates what type of memory device is used. For example, random access memory (RAM) memory can be accessed in any random order with a constant response time, thus it is typically used for cache memory and main memory. By contrast, memory device technologies that require physical movement such as magnetic disks, tapes, and optical discs, have a variable response time as the physical movement can take longer than the data transfer, thus they are typically used for secondary memory (e.g., hard drive, backup memory, etc.).\nA computer's storage system will be compliant with one or more computer storage standards that include, but are not limited to, network file system (NFS), flash file system (FFS), disk file system (DFS), small computer system interface (SCSI), internet small computer system interface (iSCSI), file transfer protocol (FTP), and web-based distributed authoring and versioning (WebDAV). These standards specify the data storage format (e.g., files, data objects, data blocks, directories, etc.) and interfacing between the computer's processing function and its storage system, which is a primary function of the computer's memory controller.\nDespite the standardization of the computer and its storage system, memory devices fail; especially commercial grade memory devices that utilize technologies incorporating physical movement (e.g., a disc drive). For example, it is fairly common for a disc drive to routinely suffer from bit level corruption and to completely fail after three years of use. One solution is to utilize a higher-grade disc drive, which adds significant cost to a computer.\nAnother solution is to utilize multiple levels of redundant disc drives to replicate the data into two or more copies. One such redundant drive approach is called redundant array of independent discs (RAID). In a RAID device, a RAID controller adds parity data to the original data before storing it across the array. The parity data is calculated from the original data such that the failure of a disc will not result in the loss of the original data. For example, RAID 5 uses three discs to protect data from the failure of a single disc. The parity data, and associated redundancy overhead data, reduces the storage capacity of three independent discs by one third (e.g., n\u22121=capacity). RAID 6 can recover from a loss of two discs and requires a minimum of four discs with a storage capacity of n\u22122.\nWhile RAID addresses the memory device failure issue, it is not without its own failure issues that affect its effectiveness, efficiency and security. For instance, as more discs are added to the array, the probability of a disc failure increases, which increases the demand for maintenance. For example, when a disc fails, it needs to be manually replaced before another disc fails and the data stored in the RAID device is lost. To reduce the risk of data loss, data on a RAID device is typically copied on to one or more other RAID devices. While this addresses the loss of data issue, it raises a security issue since multiple copies of data are available, which increases the chances of unauthorized access. Further, as the amount of data being stored grows, the overhead of RAID devices becomes a non-trivial efficiency issue."}
-{"text": "1. Field of Invention\nThis invention relates to display units such as used by retail establishments for merchandising various wares. More particularly, this invention relates to a vertically extensible bar and a clutch mechanism for holding it in adjusted position.\n2. Description of the Prior Art\nA vertically extensible bar of a display rack may have an arm for supporting a series of hangers for clothes or other merchandise. Optionally, it may cooperate with a companion bar for supporting a shelf. The typical prior art structure for holding the bar in an extended position is a series of holes in the standard and a spring detent in the bar. There are several drawbacks to this arrangement. One disadvantage is that the series of holes in the standard are unsightly. Merchandisers appreciate more elegance in the display units for their merchandise.\nOther disadvantages include the inability to achieve infinite adjusted positions, the possibility of the coupling inadvertently slipping, the necessity of performing fabrication steps both on the standard and the bar."}
-{"text": "Many articles of furniture are costly to ship because they are by nature bulky and prone to damage during transport. Therefore, it has been common to make knock down type mass-market furniture. Knock down furniture is fabricated as components, or sub-assemblies, which can be compactly packaged and economically shipped. The furniture is subsequently assembled by a retailer or a consumer using simple tools, such as common wrenches, screwdrivers, hexagonal wrenches, and the like. Most often such furniture can be subsequently disassembled, if desired. However, the advantages of knock down design will not be realized if such a design compromises the article's appearance or function, or if the article is too hard to assemble.\nWhat constitutes a compromise in appearance for a knock down article depends on an esthetic judgment, and that may vary with the individual. Nonetheless, there are some general principles which may be stated. For example, most people would conclude it is esthetically undesirable to have exposed industrial-type metal fasteners on a wooden chair. Similarly, if the knock down design involved significant changes in the proportions or shapes of the parts of a chair, compared to a traditional chair design which was obviously being emulated, then there would be a high risk that consumers would think the chair looked strange, and they would not purchase it.\nA knock down design which compromises function becomes evident when the piece of furniture is put into use. A chair may be subjected to very high loads. For instance, the chair may set on an uneven surface, a user may tilt the chair backward on the rear legs, or the chair may fall over onto a hard floor. Consequently, a knock down chair must not only have strength and rigidity when first assembled, but it must maintain such during its lifetime.\nIn furniture which is factory-assembled, it is possible to use heavy machinery and special processes. It is possible to use tight fits, diverse fasteners, and special adhesives; all to obtain the strength and durability the product demands. In contrast, by the nature of knock down furniture, there will be joints which must be made by unskilled consumers using simple hand tools. Thus, in some poorly designed knock down furniture the joints will be weak and furniture will be flimsy when initially assembled. In other such furniture, joints will loosen with time or even fail during use. In still other furniture, the knock down design may provide good strength, but be too complex for unskilled consumers to assemble correctly. And of course, a piece of knock down furniture has to be economic to manufacture, otherwise the advantage produced by lower packaging and transport costs, compared to a one-piece factory assembled chair, will be offset.\nSo, it is not easy to make a piece of knock down furniture which satisfactorily meets all the requirements. Of course, there have been many successful designs of knock down furniture. Specialized fasteners have been developed. However, certain designs of furniture by their nature still present problems which are more difficult to overcome than others. For example, joints which are made at obvious locations can be subject to inherently high stresses, as is the case when a cantilevered back rest of a chair is joined to the chair seat. Therefore, there is a continuing search for new knock down concepts and joint designs."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to a method of fabricating a Metal Oxide Semiconductor Field Effect Transistor (MOSFET), and more particularly to a method of forming a field oxide film which provides hyperfine device isolation on a Silicon-on-Insulator (SOI) substrate by means of Local Oxidation of Silicon (LOCOS).\n2. Description of the Related Art\nWith the recent remarkable progress in semiconductor devices, demand is increasing for an LSI on which both digital and analog circuits are mounted, and which performs at high speed and with reduced power consumption. To meet this demand, semiconductor devices are required to be integrated more densely. As the devices to be mounted increase in number, isolation regions must be narrower and smaller.\nA conventional method of fabricating a MOSFET in an SOI substrate by means of LOCOS is illustrated in FIGS. 2A-2F, each of which schematically shows a cross-section of the MOSFET at a fabrication step. Descriptions of the steps are as follows:\na) A pad oxide film 52 of about 5-10 nm is deposited on an SOI substrate 51. Then an active nitride film 53 of about 50-150 nm is deposited on the pad oxide film 52 as an oxidation-resistant mask (see FIG. 2A).\nb) Openings are formed in the laminated layers of the pad oxide film 52 and the active nitride film 53 at positions where field oxide films 54 are to be provided, by a conventional lithography technique (see FIG. 2B).\nc) The field oxide films 54 are formed on the SOI substrate 51 by dry oxidation (a heat treatment conducted in a dry oxygen atmosphere) (see FIG. 2C).\nd) The remaining portions of the active nitride film 53 and the pad oxide film 52 are removed (see FIG. 2D).\ne) Gate electrodes 55 are provided by a conventional process for fabricating MOSFETs (see FIG. 2E).\nf) SiO2 side walls 57 are formed by first providing an SiO2 film on the substrate and then etching back. Impurities are then introduced into the substrate by means of ion implantation to form source/drain regions 58. Finally, the impurities in the source/drain regions 58 are activated by RTA (rapid thermal annealing) and a MOSFET with low source/drain resistance is obtained (see FIG. 2F).\nIn the above-described conventional method, when the width of a field oxidation region (i.e., the distance between adjacent devices (Wi in FIG. 2B)) is reduced to 0.2 xcexcm or less (xe2x80x9csub-quarter micronxe2x80x9d), there arises a problem of insufficiency of an oxidation amount in the dry oxidation process and a resultant insufficiency in thickness of the thermal oxidation film. One of the reasons for this insufficiency in the oxidation amount is stress generated in the SOI substrate at the time of forming the openings for the field oxidation regions (in the step b).\nTo obtain a sufficient amount of oxidation, an oxidizing temperature may be increased and oxidizing time may be lengthened. However, thermal oxidation at a high temperature for a long time will cause stress in the whole SOI substrate (i.e., in the wafer). This stress may induce defects in crystals in the substrate or cause warping of the substrate. Thus, if the oxidation is conducted at high temperature for a long time to ensure a sufficient amount of oxidation in hyperfine isolation regions of about 0.2 xcexcm, the amount of oxidation will be excessively increased at areas where the design rules are less strict (e.g., peripheral circuits); i.e., the device isolation regions at those areas may be relatively wide. The thickness of the silicon layer of the SOI substrate is thinner than the conventional silicon substrate (silicon wafer). For example, the typical thickness of the silicon layer of the SOI substrate is about several nm, while the typical thickness of the conventional silicon substrate is, for example, about 625 xcexcm. Therefore, the increase of amount of oxidation may significantly cause stress in the peripheral circuit regions of the LSI, in particular, formed in the SOI, and thus cause increases in leakage currents, for example. Such effects may adversely affect the operating characteristics of the LSI which is formed on an SOI substrate.\nIn view of the aforementioned, an object of the present invention is to obtain a sufficient amount of oxidation, without changing oxidation conditions such as temperature or time, during forming of device isolation regions of 0.2 xcexcm or less by thermal oxidation.\nTo achieve the above object, a first aspect of the present invention is a method of fabricating a MOSFET, the method comprising:\n(a) preparing an SOI substrate;\n(b) depositing an oxide film on the SOI substrate;\n(c) depositing a nitride film on the oxide film;\n(d) forming an opening in the nitride film and oxide film at a predetermined region, at which a device isolation region is to be formed, by lithography for exposing a surface of the SOI substrate;\n(e) irradiating the substantially the entire area of the silicon substrate with Ar ions;\n(f) forming a field oxide film by dry oxidation; and\n(g) removing remaining portions of the nitride film and the oxide film.\nIn a second aspect of the present invention, Si ions are used in place of the Ar ions in the first aspect.\nA third aspect of the present invention is a method for fabricating a MOSFET, the method comprising:\n(a) preparing an SOI substrate having a structure of silicon layer/buried oxide/substrate;\n(b) depositing an oxide film on the SOI substrate;\n(c) depositing a nitride film on the oxide film;\n(d) forming an opening in the nitride film and oxide film at a predetermined region, at which a device isolation region is to be formed, by lithography for exposing a surface of the SOI substrate;\n(e) irradiating substantially the entire area of the SOI substrate with at least one of Ar ions and Si ions for implanting the at least one of Ar ions and Si ions into the silicon layer of the SOI substrate in the vicinity of the surface exposed by the step of forming the opening, the nitride film and the oxide film serving as a mask;\n(f) forming a field oxide film by dry oxidation; and\n(g) removing remaining portions of the nitride film and the oxide film.\nIn each aspect, the thickness of the oxide film is preferably about 5-10 nm, and the thickness of the oxidation-resistant nitride film provided on the oxide film is preferably about 50-150 nm. The ion implantation is preferably conducted at an implantation energy of 40-50 keV, and implantation dose of 1xc3x971014 to 5xc3x971015 cmxe2x88x922.\nThrough the ion implantation under these conditions, the regions of the substrate where the openings are formed become amorphous, while defects in the substrate at the regions where devices are to be mounted can be avoided. Therefore, the field oxidation is enhanced, and the thickness of the thermal oxidation film will be sufficient even at the device isolation regions having openings of 0.2 xcexcm or less. Further, no harmful effects will be caused to the electric characteristics of the device."}
-{"text": "1. Field\nMethods and apparatuses consistent with one or more exemplary embodiments relate to a method of displaying information or a user interface (UI) by a device, and the device, and more particularly, to a method of displaying appropriate information or an appropriate UI on a user device and the user device.\n2. Description of the Related Art\nWhen using various appliances such as a mobile phone, a smartphone, a laptop computer, a tablet personal computer (PC), a handheld PC, an electronic book terminal, a digital broadcasting terminal, a personal digital assistant (PDA), a portable multimedia player (PMP), a navigation device, or a smart television (TV), a user may arrange a widget or an application execution icon on a background screen or a home screen.\nHowever, according to the related art, a user background screen or a home screen of a user device is fixed regardless of information desired by a user, thus providing unwanted information or an unwanted UI to a user."}
-{"text": "Development of substances used in a variety of applications often requires an understanding of how the substances move through materials. For example, an ability of a substance (e.g., drugs, chemicals treatments, and various particulates) to diffuse through a semi-permeable material construct can provide insight into an effectiveness or a toxicity of the substance, as well as characteristics of the material construct. In some implementations, diffusion cells can be used to examine such parameters."}
-{"text": "Most electronic equipment, and in particular computers, utilize a series of chips which are connected to a motherboard in order to form the signal processing part of the equipment. Various chips may assume a single function or multiple functions which are used by the equipment. The group of chips used together is sometimes referred to as a chipset.\nFIG. 1 is a block diagram showing the arrangement of a chipset on a motherboard for a computer. The chip set 100 includes a first chip 102 which carries the central processing unit for the device. Memory controller hub 104 acts as a central controller to move data into and out of memory and to other related chips. Chip 106 is a graphics chip which generates various graphic arrangements for display. Chip 108 is the memory itself, either RAM or ROM memory. Chip 110 is an input/output controller hub which transfers data to various input/output devices. Chip 112 includes connections to a hard disk drive. Chip 114 is a chip which connects to other peripheral components.\nTypically, each chip in a chip set is formed of two parts. The first part is the core which is the circuitry which handles the main function of the device itself. Also on the chip are input/output circuits for connecting the core to other chips. For example, the memory controller 104 would have a central core and an input/output device connected to each of the four other chips 102, 106, 108 and 110 to which it is connected.\nFor every pair of chips that are connected, an interface is provided to connect the input/output devices of the chips to each other. Thus, the CPU 102 and memory controller hub 104 are connected by a front side bus (FSB) 116. Likewise, memory controller hub 104 is connected to graphics chip 106 through the advanced graphics port (AGP) 118. Memory 108 is connected to the memory controller hub 104 by a system memory bus 120. Memory controller hub 104 is connected to the input/output controller hub 110 through hub link 122. The input/output controller hub 110 is connected to the hard disk drive 112 through IDE 124. The I/O controller hub 110 is connected to the peripheral components chip 114 through the peripheral components interface 126.\nFIG. 1 also shows a clock circuit 113 which is another chip connected on the motherboard. This clock provides clock signals of various frequencies to the various other chips. These particular connections are not specifically shown but all chips on the motherboard are connected thereto to receive clock signals which are necessary for the synchronization of the entire device.\nSome of the interfaces on the motherboard are considered to be source synchronous interfaces. In the present example, the front side bus 116, the advanced graphics port 118 and the hub link 122 are all source synchronous circuits. On the other hand, a system memory bus 120 and IDE 124 are not source synchronous interfaces. In such an interface, data signals and strobe signals are used to transfer data in a synchronous fashion. These signals occur in a certain preset timing relationship so that data being transferred can be expected at a particular time location."}
-{"text": "1. Technical Field\nThe present invention relates to a switch device suitable for starting a vehicle engine.\n2. Related Art\nRecently a type of vehicle, in which a user does not conventionally insert a key in a key cylinder to turn the key, but the user having a proper electronic key starts up an engine only by pressing a push button of an engine starting switch device provided on a driver seat on a condition that the vehicle is equipped with an authentication system such as a so-called immobilizer, has become widespread in vehicles such as a four-wheeled vehicle. Japanese Unexamined Patent Publication No. 10-205183 discloses an automotive key cylinder in which a drain property is considered. In the automotive key cylinder disclosed in Japanese Unexamined Patent Publication No. 10-205183, a drain hole is made in a lower portion on a front-end side of a case, and a liquid (such as rain water) invading in a cylinder head from a key plate hole is drained away from the drainage hole to the outside of the case."}
-{"text": "1. Field of the Invention\nThe invention relates to systems used for chemical sterilization of medical devices, and more particularly, to systems having multiple chambers used for chemical sterilization of medical devices.\n2. Description of the Related Art\nMedical instruments have traditionally been sterilized using either heat, such as is provided by steam, or a chemical, in the gas or vapor state. Sterilization using hydrogen peroxide vapor has been shown to have some advantages over other chemical sterilization processes.\nThe combination of hydrogen peroxide with a plasma provides certain additional advantages, as disclosed in U.S. Pat. No. 4,643,876, issued Feb. 17, 1987 to Jacobs et al. U.S. Pat. No. 4,756,882, issued Jul. 12, 1988 also to Jacobs et al. discloses the use of hydrogen peroxide vapor, generated from an aqueous solution of hydrogen peroxide, as a precursor of the reactive species generated by a plasma generator. The combination of hydrogen peroxide vapor diffusing into close proximity with the article to be sterilized and plasma acts to sterilize the articles and remove residual hydrogen peroxide. However, effective sterilization of articles having long narrow lumens are very difficult to achieve, since the methods are dependent upon diffusion of the sterilant vapor into close proximity with the article before sterilization can be achieved. Thus, these methods have been found to require high concentration of sterilant, extended exposure time and/or elevated temperatures when used on long, narrow lumens. For example, lumens longer than 27 cm and/or having an internal diameter of less than 0.3 cm have been particularly difficult to sterilize. The sterilization of articles containing long narrow lumens therefore presents a special challenge.\nU.S. Pat. No. 4,744,951 to Cummings et al. discloses a two-chambered system which provides hydrogen peroxide in vapor form for use in sterilization processes. The sterilant is initially vaporized in one chamber and then applied to the object to be sanitized in another single sterilizing chamber, thereby producing a concentrated hydrogen peroxide vapor which is relatively more effective. The sterilization processes are designed for furnishing concentrated hydrogen peroxide vapor to interior surfaces of articles having a tortuous or a narrow path. However, the sterilization processes are ineffective at rapidly sterilizing lumened devices, since they depend on the diffusion of the hydrogen peroxide vapor into the lumen to effect sterilization.\nU.S. Pat. No. 4,797,255 to Hatanaka et al. discloses a two-chambered sterilization and filling system consisting of a single sterilization chamber adjacent to a germ-free chamber utilized for drying and filling sterilized containers.\nU.S. Pat. No. 4,863,688 to Schmidt et al. discloses a sterilization system consisting of a liquid hydrogen peroxide vaporization chamber and an enclosure for sterilization. The enclosure additionally may hold containers wherein the hydrogen peroxide sterilant vapor does not contact the interior of the containers. This system is designed for controlling the exposure to the hydrogen peroxide vapor. The system is not designed for sterilizing a lumen device.\nU.S. Pat. No. 4,952,370 to Cummings et al. discloses a sterilization process wherein aqueous hydrogen peroxide vapor is first condensed on the article to be sterilized, and then a source of vacuum is applied to the sterilization chamber to evaporate the water and hydrogen peroxide from the article. This method is suitable to sterilize surfaces, however, it is ineffective at rapidly sterilizing lumened devices, since it too depends on the diffusion of the hydrogen peroxide vapor into the lumen to effect sterilization.\nU.S. Pat. No. 4,943,414, entitled \u201cMethod for Vapor Sterilization of Articles Having Lumens,\u201d and issued to Jacobs et al., discloses a process in which a vessel containing a small amount of a vaporizable liquid sterilant solution is attached to a lumen, and the sterilant vaporizes and flows directly into the lumen of the article as the pressure is reduced during the sterilization cycle. This system has the advantage that the water and hydrogen peroxide vapor are pulled through the lumen by the pressure differential that exists, increasing the sterilization rate for lumens, but it has the disadvantage that the vessel needs to be attached to each lumen to be sterilized.\nU.S. Pat. Nos. 4,937,046, 5,118,471 and 5,227,132 to Anderson et al. each disclose a sterilization system which uses ethylene oxide gas for sanitation purposes. The gas is initially in a small first enclosure and thereafter slowly permeates into a second enclosure where the objects to be sterilized are located. A medium is then introduced into the second enclosure to flush out the sterilizing gas into a third enclosure containing the second enclosure. An exhaust system then exhausts the sterilant gas and air from the third enclosure. These systems also have the disadvantage of relying on the diffusion of the sterilant vapor to effect sterilization and hence are not suitable for rapidly sterilizing lumened devices.\nU.S. Pat. No. 5,122,344 to Schmoegner discloses a chemical sterilizer system for sterilizing items by vaporizing a liquid chemical sterilant in a sterilizing chamber. Pre-evacuation of the sterilizer chamber enhances the sterilizing activity. Sterilant is injected into the sterilizer chamber from a second prefilled shot chamber. This system also relies upon diffusion of sterilant vapor to effect sterilization and is also not suitable for rapidly sterilizing lumened devices.\nU.S. Pat. No. 5,266,275 to Faddis discloses a sterilization system for disinfecting instruments. The sterilization system contains a primary sterilization chamber and a secondary safety chamber. The secondary safety chamber provides for sensing and venting to a destruction chamber any sterilization agent that is released from the primary sterilization chamber. This system, as in other systems, also relies upon diffusion of sterilant vapor to effect sterilization and is also not suitable for rapidly sterilizing lumened devices.\nIn U.S. Pat. Nos. 5,492,672 and 5,556,607 to Childers et al, there is disclosed a process and apparatus respectively for sterilizing narrow lumens. This process and apparatus uses a multicomponent sterilant vapor and requires successive alternating periods of flow of sterilant vapor and discontinuance of such flow. A complex apparatus is used to accomplish the method. Additionally, the process and apparatus of '672 and '607 require maintaining the pressure in the sterilization chamber at a predetermined subatmospheric pressure.\nIn U.S. Pat. No. 5,527,508 to Childers et al., a method of enhancing the penetration of low vapor pressure chemical vapor sterilants into the apertures and openings of complex objects is disclosed. The method repeatedly introduces air or an inert gas into the closed sterilization chamber in an amount effective to raise the pressure to a subatmospheric pressure to drive the diffused sterilant vapor further into the article to achieve sterilization. The '508, '672 and '607 Childers inventions are similar in that all three require repeated pulsations of sterilant vapor flow and maintenance of the sterilization chamber pressure at a predetermined subatmospheric pressure.\nIn U.S. Pat. No. 5,534,221 to Hillebrenner et al., a device and method for sterilizing and storing an endoscope or other lumened medical device is disclosed. The device includes a sealable cassette in which the endoscope or other medical device is placed. The cassette has an input port for receiving a sterilizing agent through a connector, an output port for expelling the sterilizing agent when a vacuum is applied thereto through a connector, and check valves in the input and output ports to open the ports when the connectors are coupled to the ports and to seal the ports when the connectors are removed from the ports such that after the endoscope has been sterilized, it remains sterilized within the cassette until the cassette is opened. The method of the '221 invention involves placing the medical device inside the cassette and coupling the device to either the input or output port of the cassette. The cassette is then placed in an outer oven-like container or warming chamber where the temperature is properly maintained. Connections are made to open the input and output ports on the cassette such that the sterilizing agent may be introduced through a first port to bathe the outside of the medical instrument or other object, such as an endoscope while one end of the hollow object, such as the endoscope, is coupled to the output port where a vacuum is supplied external to the cassette to pull the sterilization agent into the cassette and through the interior passageways of the endoscope. When the sterilization process is completed, the warming chamber is opened and the sterilizing cassette is simply removed from the chamber with the input and output ports being uncoupled from their respective sources. A tight seal is maintained and the object remains in the sterilized interior of the cassette until the cassette is opened or the device is to be used. Thus, the '221 invention is concerned with providing a means whereby a sterilized medical device can be retained within a cassette in which it was sterilized until ready for use, thus avoiding any contamination by exposure to the atmosphere or handling before use. Additionally, in some cases of the '221 invention, wherein the lumen of the device to be sterilized is connected to the output port, particularly wherein the devices have long, narrow lumens, the time to expel the sterilizing agent through the lumen and out of the cassette may be undesirably long. Also, in cases wherein the lumen device is very flexible, lumen collapse may occur, either slowing or preventing vapor exit or causing lumen damage.\nU.S. Pat. Nos. 5,445,792 and 5,508,009 to Rickloff et al. each disclose a sterilization system essentially equivalent to the system disclosed in Hillebrenner '221.\nU.S. Pat. No. 5,443,801 to Langford teaches a transportable cleaning/sterilizing apparatus and a method for inside-outside sterilization of medical/dental instruments. The apparatus avoids the use of heat, pressure, severe agitation, or corrosive chemicals which might damage delicate equipment. This invention uses ozone gas or solution as sterilant. It does not involve the use of sterilant vapor or vaporizing a sterilant solution into vapor, and is not suitable for operations under vacuum because flexible bags or containers are used.\nIn consideration of the foregoing, no simple, safe, effective method of sterilizing smaller lumens exists in the prior art. Thus, there remains a need for a simple and effective method of vapor sterilization of articles with both long, narrow lumens as well as shorter, wider lumens. Furthermore, there also remains a need for a simple and effective sterilization system with independently operable chambers."}
-{"text": "A mass spectrometer is an instrument used to measure the mass, or more specifically the mass to charge ratio, of ionized atoms or electrically charged particles. Mass spectrometers help determine the composition of an unknown sample by isolating ionized atoms based on their mass-to-charge ratio, measured in Atomic Mass Units per charge (AMU/q). Mass spectrometers find widespread application in the basic sciences, medicine, and space-based research. Two common space-related applications of mass spectrometry are the study of the composition of planetary atmospheres and the monitoring of air quality on manned space missions. Although mass spectrometry has been used in space-related applications for many years, usage in space presents unique design challenges, both in terms of detection sensitivity and logistical considerations such as weight and power requirements.\nIn early mass spectrometers, atoms or molecules were ionized by a hot filament and accelerated through the instrument under the influence of voltage gradients. The ions followed a semi-circular trajectory through the instrument, which utilized a strong magnetic field to selectively direct ions of a specific mass towards a detector. By controlling the strength of the magnetic field and the accelerating voltage, ions of different mass/charge ratios could be selectively guided towards the detector. These early mass spectrometers suffered from numerous deficiencies and drawbacks, most significantly the difficulty in achieving and maintaining a stable magnetic field.\nQuadrupole mass spectrometers (QMS) eliminated the need for magnetic fields. Similar to its predecessor, a QMS employs a hot filament to ionize the atoms or molecules. Ionization results from the conversion of normally neutral atoms or molecules to electrically charged particles. The ions are accelerated through a mass filter having four parallel metal rods, referred to as the quadrupole. DC and RF (frequency \u03a9) voltages are applied to opposing pairs of these rods with opposite polarities to create an electric field inside the rod assembly. For a given DC and RF voltage, only ions of a certain mass-to-charge ratio will pass through the quadrupole filter, while all other ions are thrown out of their original path. The stability region is defined by ion trajectories that are periodic and bounded. A detector placed at the end of the rod assembly opposite the ionizer measures those ions that pass through the quadrupole filter. A mass spectrum is obtained at the detector by measuring the ions passing through the quadrupole filter as the voltages on the quadrupole rods are varied. The mass resolution of the QMS is the maximum atomic mass that can be distinguished. The maximum attainable resolution is determined by both the fidelity of the electronics and the overall tolerances of the instrument design. Generally, the voltages employed in QMS systems are of the order of a few thousand volts to obtain a mass resolution of a few hundred Daltons.\nThe quadrupole rods can be a circular or hyperbolic. Circular rods are easier to manufacture and consequently cheaper. However, the quadrupole electric field produced with circular rods is slightly distorted, which can reduce the maximum attainable mass resolution of the instrument. Consequently, in applications requiring high mass resolution, the more difficult and expensive to manufacture hyperbolic rods are employed as quadrupole rods.\nTo improve resolution, the electric field generated by the quadrupole rods can be perturbed by introducing an excitation RF signal at an auxiliary frequency (\u03c9) different from the fundamental frequency (\u03a9). This perturbation causes the original stability region to break into smaller regions termed islands, including an \u2018upper stability island.\u2019 The result of this auxiliary frequency is the creation of bands of instability in the previously stable regions of the electric field. Charged particles having a mass within a certain range that previously passed through a stable region of the electric field may now be thrown off trajectory as they coincide with these islands of instability. In this way, the use of an auxiliary frequency to drive the quadrupole rods allows a QMS to operate with improved resolution. A QMS driven under an auxiliary frequency excitation is able to better differentiate between charged particles having close, yet different masses, or mass-to-charge ratios. The size and shape of the upper stability island is determined by the auxiliary frequency used and the amplitude of the excitation RF signal. To create an island of appropriate size, for example, the auxiliary frequency (\u03c9) inserted into the QMS system needs to be near an integer multiple of the fundamental RF frequency (i.e., \u03c9=0-0.1\u03a9, 0.9-1.1\u03a9, 1.9-2.1\u03a9, etc.). Employing an excitation RF signal in one of these auxiliary frequency ranges in conjunction with the DC voltage U and the RF voltage V allows for improved resolution and discrimination between ions with small differences in their mass-to-charge ratio. In general, the auxiliary signal amplitude required for appropriate island formation increases with auxiliary signal frequency.\nUnfortunately, this use of auxiliary frequency excitation presents problems in constrained applications, such as space-based applications, where it is advantageous for the QMS to have increased sensitivity and enhanced resolution to better detect and differentiate between complex molecules with higher masses. Having to excite the quadrupoles with an excitation RF signal at an auxiliary frequency in order to create islands of stability/instability requires higher power and increased complexity of the voltage control system. Additionally, the excitation RF signal must be driven at an amplitude that corresponds to a few hundred volts (\u02dc10% of the fundamental RF signal amplitude), to create islands of the appropriate size. However, in space-based applications, power and size is at a premium.\nOther factors that affect the resolution and accuracy of the measurement made by the QMS are imperfections in the rods and limitations of the electronics. Furthermore, electronic component values drift with temperature and time, which can have the material effect of shifting the operating point of the quadrupole sufficiently to degrade the detected mass spectrum.\nThus there exists a need to enable a QMS to resolve species of heavy, complex molecules in a power efficient manner, but also to improve the tolerance of the instrument to variations in electronic component values."}
-{"text": "A virtual machine (VM) is a software implementation of a physical computer. Computer programs designed to execute on the physical machine execute in a similar way when executed on a VM. A VM provides a complete system platform to support a full operating system (OS). A physical machine can be shared between users by using different VMs, each running a different OS.\nModern processor architectures have enabled virtualization techniques that allow multiple operating systems and VMs to run on a single physical machine. These techniques use a hypervisor layer that runs directly on the physical hardware and mediates accesses to physical hardware by providing a virtual hardware layer to the operating systems running in each virtual machine. The hypervisor can operate on the physical machine in conjunction with a \u2018native VM\u2019. Alternatively, the hypervisor can operate within an operating system running on the physical machine, in conjunction with a \u2018hosted VM\u2019 operating at a higher software level.\nExamples of VM technology are:                Linux Kernel-Based Virtual Machine (KVM) allows one or more Linux or Windows virtual machines to be run on top of an underlying Linux that runs KVM.        Xen allows a guest (virtualized) Linux to be run on top of Linux.        Parallels allows Linux and Windows on top of Mac OS X.        VMWare allows Linux and Windows systems on top of Mac OS X, Windows and Linux systems.        \nA user may want to migrate a workload operating on one physical machine (host A) to another physical machine (host B), for example, for machine maintenance or for performance optimisation. If the instruction set architecture is the same on both host A and host B, the VM needs to be shut down on host A, restarted on host B, and the workload migrated. However, if the ISA on each physical machine is different, a migration is problematic, because, for example, the format state of the VM on host A is inappropriate for the format state of the VM on host B."}
-{"text": "1. Field of the Invention\nThe present invention relates to a face image obtaining apparatus for obtaining a face image to be attached to a personal paper or the like belonging to an individual, in particular, ID card, magnetic card, or the like. More specifically, the present invention relates to a face image obtaining apparatus having hand-related biological information obtaining function, as well as providing a face image of the user.\n2. Description of the Related Art\nCurrently, unattended face image obtaining apparatuses for providing face images of the users are installed on the street. Such apparatuses provide recorded face images to the users by printing on plain papers or stickers.\nIn addition, a face image obtaining apparatus capable of obtaining biological information of the user such as fingerprint and the like, as well as face image, is proposed as described, for example, in International Patent Publication No. WO2005/050508. In the apparatus disclosed in the aforementioned patent publication, a face image of the user is recorded first, then the biological information. The apparatus records monitoring photographs, including a face image of the user, before and after obtaining biological information in order to provide a proof record when the user is switched for counterfeiting purpose. Such monitoring photographs may have a deterrent effect on the counterfeiting user switching. But, it is difficult to prevent such counterfeiting user switching at the site where the biological information and face image are obtained.\nIn order to prevent such counterfeiting user switching, a face image obtaining apparatus in which the face image and biological information are obtained at the same time has also been considered. If, for example, a fingerprint is obtained as the biological information, it is not an easy task for the user not accustomed to taking a fingerprint to obtain a face image and a proper fingerprint applicable to fingerprint authentication at the same time. When a user initially failed to obtain a fingerprint, if reacquisition of the fingerprint is authorized, the counterfeiting user switching may not be prevented.\nIn the mean time, another face image obtaining apparatus is also proposed as described, for example, in International Patent Publication No. WO2005/050508. In the apparatus, biological information is obtained before a face image, and when obtaining the face image, the biological information is obtained again to verify the identity of the user.\nThe face image obtaining means described in Japanese Unexamined Patent Publication No. 2005-141429 may prevent the counterfeiting user switching, but has a problem that the biological information needs to be reacquired when recording the face image, which increases the burden on the user.\nThe present invention has been developed in view of the problem described above, and it is an object of the present invention to provide a face image obtaining apparatus capable of reliably obtaining biological information applicable to authentication, and preventing counterfeiting user switching without increasing the burden on the user."}
-{"text": "1. Field of the Invention\nThe present invention relates to storage managers that provide data storage services to software applications. More particularly, the invention concerns the provision of filtering functions such as encryption, compression and other data conversions as part of storage manager operations.\n2. Description of the Prior Art\nBy way of background, a storage manager is a system that acts as an intermediary between a software application (such as a backup/restore program or a web server) and a data storage resource (such as a tape drive, a disk drive, a storage subsystem, etc.). The storage manager, which could be integrated with the application program or implemented separately therefrom, provides an interface that accepts objects for storage and subsequently retrieves the objects upon request. Applications for which a storage manager has been used include the management of backup images of database installations, enterprise application data, individual workstations, web content, etc.\nThere is often a need for a storage manager to filter the data being written to or read from physical storage devices by compressing or encrypting the data. Existing storage managers that provide support for compression and/or encryption do so in one of two ways. Most commonly, such filtering is provided by algorithms that are embedded in the storage manager product itself. Less commonly, such filtering is supported by providing a programming hook that gives a storage manager user the option of writing their own algorithm(s). With this option, the user is also required to re-implement much of the functionality of the storage manager on their own.\nDrawbacks of the first approach include:                The user is limited to the compression and/or encryption algorithms that are built into the storage manager product.        Some products support encryption but not compression and vice versa.        Some products support only weak encryption or poor compression.        The storage manager vendor may charge customers extra to enable the compression and/or encryption algorithms that are built in.        If a built-in algorithm is found to have a security flaw or a crippling bug, a customer cannot immediately swap in a different off-the-shelf algorithm to avoid exposure to the risk.        Storage manager customers must wait for the vendor to update the embedded algorithms with the latest technology when better algorithms are invented, even though the new technology may already exist in stand-alone off-the-shelf programs.        A vendor may not implement a particular compression or encryption algorithm that a customer desires.        \nDrawbacks of the second approach include:                The storage manager programming hook places a burden on the customer to re-implement much of the functionality the storage manager otherwise provides. The user must typically write a program that can accept objects for storage, track the location of these objects, write and read them to/from physical storage devices, and retrieve them upon request based on whatever query protocol the storage manager requires, as well as write in the desired compression and/or encryption algorithms. In this solution, the storage manager essentially delegates all work to the user and does not provide any functionality of its own. The storage manager mostly acts as a hollow shell or \u201cstub\u201d that forwards all storage and retrieval requests to the user-written external program for handling. The storage manager itself merely assembles and disassembles buffers of information that pass between it and the application that is calling it, and provides stubs for the interface APIs (Application Programming Interfaces) but delegates most of the work to the user's program.        This approach provides very little support for compression and encryption. There is the programming hook but the customers are required to create the needed support at great additional expense and effort to themselves.        A customer who uses the programming hook but does not sufficiently test and debug their external program may find that their data has been corrupted by their own custom program, or that bugs in the program prevent the retrieval of storage objects at a critical time, such as when they need them to restore a down system.        If the event described in the preceding paragraph occurs, the storage manager vendor may find itself exposed to liability for the customer's own programming mistakes.        \nAccordingly, a need exists for a storage manager filtering technique that overcomes the foregoing disadvantages. What is required is a solution that allows storage manager filters to be easily implemented without having to redesign the storage manager or duplicate its functionality in a custom program. It would be further desirable to provide the capability of implementing new and different filters. At present, the most common needs for storage manager filtering are compression and encryption. However, it is submitted that the possibilities are broader, and it may be advantageous in some circumstances to provide other data conversions, such as converting between English and metric units, or between different code pages or character sets like ASCII (American Standard Code for Information Interchange), EBCDIC (Extended Binary Coded Decimal Interchange Code), and Unicode. By way of example, this capability would be useful if backup data was generated by a first system in a first character format (e.g., a mainframe computer using EBCDIC character) and the data needed to be restored to a second system that used a second character format (e.g., a workstation using ASCII character encoding). Another area where storage manager filtering could be used is the generation of audit trails. Such a filter could be used to inspect the data being stored or retrieved and generate audit information for management purposes."}
-{"text": "Since its introduction in 1975, the well-known Kohler and Milstein technique (Nature 256:495, 1975) for the production of mouse hybridoma cells has made it possible to produce large quantities of mouse antigen-specific monoclonal antibodies that are useful in a number of investigative, diagnostic and therapeutic applications. The mouse hybridoma cells, which are initially produced by the fusion of antibody-producing cells (B-lymphocyte cells, hereinafter referred to as B-cells) with malignant, transformed B-cells (in vivo transformed, myeloma cells from mice afflicted with myeloma or plasmacytoma) are capable of producing large quantities of monoclonal antibodies with predetermined specificities.\nUsing the Kohler and Milstein technique, a B-cell and a plasmacytoma cell are fused using, for instance, polyethylene glycol, lysolecithin or Sendai virus as the cell-fusing agents. A selectable marker must be present in the fused cells to enable them to be selected from parent cells and other non-hybridoma cells. As an example, the plasmacytoma fusion partner is generally deficient in an enzyme, (for instance, hypoxanthine-guanosyl phosphoribotransferase (HGPRT)) that is necessary for growth of the fused cell in certain media (hypoxanthine-, aminoprotein-, and thymidine- containing medium or HAT medium). This enzyme deficiency enables the resultant hybrids to be selected for their ability to grow in such media. This insures that only B-cell:plasmacytoma cell hybrids are recovered since neither parental cells (nor hybrids comprising B-cell: B-cell and plasmacytoma: plasmacytoma cell) can survive in selective media.\nMurine antibodies produced with the Kohler and Milstein technique are generally unsuitable for administration to human subjects as in-vivo therapeutic agents, e.g., to provide passive immunity to an infectious agent. The extension of the Kohler and Milstein hybridoma technology to the production of human monoclonal antibodies has been limited, largely due to: (1) the lack of good human plasmacytoma cells for use as fusion partners; (2) the low frequency of cell fusion events (\"fusion efficiency\"); and (3) the relative scarcity of B-cells circulating in human blood and producing specific antibodies against antigens of interest (and the inherent difficulties in isolating such cells). These factors make it difficult to obtain hybridoma cell lines secreting human monoclonal antibodies of a predetermined specificity.\nCasali et al. (Science 234:476-479, 1986) disclosed a method which represents some progress toward making human monoclonal antibody-producing cells. Normal B-cells obtained from peripheral human blood were pre-selected for their specificity to a given antigen using Fluorescence-Activated Cell Sorting (FACS). Positively selected clones were then established as lymphoblastoid cells in vitro by infecting such cells with Epstein-Barr virus (EBV). The EBV infected cells produced antigen-specific human monoclonal antibodies. However, the method of Casali et al. has several significant drawbacks which impair its usefulness: (1) the amount of monoclonal antibodies produced by the Casali et al. cells is relatively low, and (2) the antibody producing cells are relatively unstable and some clones stop antibody production prematurely. In addition maintenance of the antigen-specific antibody production requires repeated cloning of the cells, a time-consuming and inefficient procedure given the low clonogenic (i.e. growth) properties of the resultant lymphoblastoid or lymphoblastoid cell lines (LCL); (3) large-scale production and purification of the monoclonal antibodies is inefficient in view of the long doubling time and high serum requirements of the LCL; and (4) the LCL produced by this process cannot be grown as tumors in animals. Such tumor cell growth permits the amplification and purification of antibodies from ascitic fluids, an efficient method for large scale antibody production that is widely used in making murine monoclonal antibodies. Finally, the Casali et al. method does not dispense with the requirement for identifying a human B-cell specific to a certain antigen.\nCopending U.S. patent application Ser. No. 041,803 (allowed) filed Apr. 23, 1987 of Riccardo Dalla-Favera discloses a method for the production of human monoclonal antibody-producing cells. Specific B lymphocytes are selected using the method of Casali et al. (supra), infected with Epstein Barr virus (EBV) and transfected with activated c-myc DNA sequences. The resultant cells are tumorigenic (i.e. can grow in semisolid medium and animals such as rats or mice) and clonogenic and produce monoclonal antibodies of a predetermined specificity. However, it was found that these cells still produce relatively low amounts of antibody because the transfected lymphoblastoid cells had not undergone differentiation.\nCurrently there is no convenient and reliable system available for the production of human monoclonal antibodies wherein the monoclonal antibody-producing cells are stable, highly malignant and which can be readily manipulated to produce high antibody titers.\nIt, is therefore an object of the present invention to provide a method for the production of tumorigenic human cells that are capable of producing human monoclonal antibodies.\nA further object of the present invention is to provide a transformed lymphoblastoid cell that is useful as a fusion partner in the production of human monoclonal antibodies.\nAnother object of the present invention is to provide a transformed lymphoblastoid cell that demonstrates high level proliferative, differentiation and antibody production properties.\nAnother object of the present invention is to produce a new human cell line comprising human B-cells infected with Epstein-Barr virus and which have at least one exogenous activated K-, N- or H-ras oncogene DNA sequences."}
-{"text": "The invention concerns cosmetic compositions for the treatment of hair or skin, having a content of new, macromolecular compounds derived from chitosan, which are employed in a suitable cosmetic foundation.\nThe invention further concerns new N-hydroxypropyl-chitosans, as well as processes for the production thereof.\nIt is already known to employ cationic polymers, in particular polymers which display quaternary ammonium groups, as conditioning agent in cosmetic compositions, particularly for the treatment of hair. Based upon a reciprocal action between their ammonium groups and the anionic groups of the hair, the cationic polymers possess a great affinity towards keratin fibers.\nIt has been determined that the employment of such cation-active polymers in such cosmetic compositions provides numerous advantages: the disentanglement of the hair, as well as its treatment, are facilitated, and, moreover, the hair is provided with elasticity and lustrous effect. On account of their affinity towards keratin, however, these polymers tend to accumulate in the hair upon repeated use, so that the hair becomes heavier, which is undesirable as a final effect.\nMoreover, synthetic polymers provide problems on account of the physiological activity of possibly present trace monomers, which are removable from the polymers only with difficulty.\nIt has already been attempted to eliminate the above-mentioned disadvantages by emplying in such cosmetic compositions the water-soluble salts of chitosan, i.e. polyglucosamines producable from chitin by means of entacetylation. In this connection, reference is made to European Patent 0 002 506, as well as German Pat. No. 26 27 419.\nIn the same manner as with the plurality of cation-active polymers having quaternary groupings, chitosan likewise frequently provides the disadvantage that it is not too compatible with the anion-active surface-active agents which in customary manner find use in cosmetic compositions for the treatment of hair, particularly in shampoos. It is therefore necessary to apply the chitosan for penetration in separate treatments, namely before and/or after the shampooing.\nIn addition, the chitosan displays, in neutral and alkaline medium, near insolubility, so that its use, for example, in alkaline permanent shaping compositions or hair dyeing compositions, is not possible.\nBy means of employment of glycidyl chitosans instead of chitosan salts according to DE-OS 32 No. 23 423, the above-mentioned disadvantages can be avoided. The reaction of chitosan with glycide is, however, very cost-intensive, since glycide is a more expensive raw material, not produced on a large scale."}
-{"text": "The invention relates to handles for carrying batteries, and more particularly to a rope-type battery carrying handle that has an end of the rope removably attached to an end of a grip.\nStarting, lighting, and ignition (SLI) batteries are typically used in automotive, recreational, and other applications, are heavy, cumbersome, and usually require two hands, or often two people, for carrying. The desirability of providing such batteries with attachable/detachable handles for facilitating carrying, placement, and retrieval of such batteries has long been known. Such handles are a particular convenience in batteries designed for use in boats or in uninterrupted power supply (UPS) applications which must be frequently moved for storage, service, or recharging.\nBail-type handles, which are known in the art, typically comprise a U-shaped or C-shaped member attached to opposing sides of a battery casing, either on its container or cover. With such handles, the battery may be carried in much the same fashion as a picnic basket or bail.\nSubstantially rigid bail-type handles are known in the art. A variety of such handle designs have been proposed for carrying batteries. Detachable, substantially rigid bail handles are disclosed, for example, in U.S. Pat. No. 3,093,515 to Rector, U.S. Pat. No. 3,956,022 to Fox, U.S. Pat. No.4,029,248 to Lee, U.S. Patent No. 4,673,625 to McCartney et al., U.S. Patent No. 5,232,796 to Baumgartner, U.S. Pat. No. 5,242,769 to Cole et al., and U.S. Pat. Des. No. 292,696 to Sahli.\nRope-type handles are likewise known in the art. Rope-type handles typically have one or more injection molded plastic part coupled by flexible rope sections and, accordingly, are physically highly flexible. The rope sections are generally a braided synthetic material such as polypropylene.\nAccording to one type of rope handle design, the ends of the rope handle are manually fed into two holes and coupled to the battery container. In the battery disclosed in U.S. Pat. No. 3,092,520 to Buskirk et al., the rope handle is coupled to the battery container by cementing the ends of the rope in recesses in projections on the sides of the battery container. Alternately, the ends of the rope handle may include an enlarged molded plastic portion and may be pressed into slots underneath the handle bracket area as shown, for example, in U.S. Pat. No. 3,797,876 to Gummelt and U.S. Pat. No. 4,013,819 to Grabb. According to other designs, the ends of the rope may be enlarged as shown for example in British Patent 869,329, or the ends coupled or welded together as shown for example in British Patent 869,329 and British Patent 1,453,977.\nIn another type of rope handle design, looped rope portions extend from the ends of a molded plastic grip portion as shown, for example, in U.S. Pat. No. 971,876 to Apple, U.S. Pat. No. 4,791,702 to McVey, and U.S. Pat. No. 5,242,769 to Cole et al. The looped rope portions are then coupled to the battery container via dedicated protrusions extending from the walls of the battery by looping the rope around the protrusion and then securing it into a recess or the like.\nAnother such rope handle design is disclosed in U.S. Pat. No. 5,144,719 to Arthur. The Arthur patent discloses a xe2x80x9cU-shapedxe2x80x9d handle having one end of the rope embedded in one depending leg of the handle. The opposite end of the rope includes an enlarged head, which may be fed through lugs on the battery. The enlarged head and the adjacent length of rope are then laid into a tri-part vertical slot on the other depending leg of the handle, the head being disposed in the upper portion of the slot, the adjacent rope extending through the lower two portions of the slot. Significantly, however, the head and adjacent rope section are not secured to the handle. As may be seen in the illustrations of the reference, there is sufficient clearance between the head and the slot, as well as the adjacent rope section and the slot such that the head and rope section may become easily dislodged from the handle leg unless a constant vertical force is maintained on the handle. Accordingly, the Arthur handle does not provide an attachment mechanism which is reliable. Moreover, the intricate coupling design requires the user to have a high level of manual dexterity and a working knowledge of the defailed structure of the complex attachment.\nInstallation of these rope handle designs may be labor intensive. Properly securing the ends of the rope to the battery container or securing the loop ends around a protrusion and into a recess can be quite time consuming and may require manual dexterity. These difficulties in installing the battery handles can lead to improper installation, which can result in an unreliable battery handle.\nAdditionally, these designs generally require specialized handle brackets to be molded into specific containers. Complicated grip and/or rope end configurations may also be required. These requirements can result in increased costs in the form of mold and tooling costs, as well as increased labor and downtime costs during changeover. Further, storage and floor space costs increase because the battery manufacturer must maintain larger inventories.\nIt is a primary object of the invention to provide a rope handle that may be reliably and easily assembled onto a battery container and which remains securely coupled to the battery until purposely removed by the user.\nA related object of the invention is to provide a rope handle arrangement that has a relatively simple design, and does not require high manual dexterity to assembly for a secure, reliable handle.\nIt is a further object of the invention to provide a rope handle that may be utilized with a battery that produces an acceptable appearance.\nIt is another object of the invention to provide a rope handle that contributes to the production of an economical battery. A related object of the invention is to provide a rope handle design that minimizes manufacturing and inventory costs.\nThese and other objects and advantages of the invention will be apparent to those skilled in the art upon reading the following summary and detailed description and upon reference to the drawings.\nIn accomplishing these and other objects of the invention, there is provided a battery that includes rope handles each of which engages a handle bracket on an end wall of a conventional battery container. Each rope handle includes a grip with a retaining recess at one end of the grip and a rope secured to the other end of the grip by molding or the like. The rope has an enlarged distal end or a cylindrical plug molded for engaging the retaining recess of the grip. The retaining recess includes a generally keyhole-shaped slot which extends through the grip from a first surface to a second surface and which has a hole portion and a channel portion projecting radially from the hole portion and terminating at an end. The retaining recess also includes a counterbore located on the first surface of the grip and encompassing the hole portion of the slot. To secure the rope to the grip, the rope is slid through the slot and the plug is subsequently drawn towards the grip and is retained within the counterbore, thus securing the handle to the battery container. In other words, the retaining recess includes a counterbore with a subjacent retaining surface for receiving and supporting the plug, and radially extending slot. The rope is laterally advanced through the slot to move the plug into position above the counterbore. The plug is then pushed down into position in the counterbore and/or a downward force is exerted on the rope to position the plug and secure the rope handle."}
-{"text": "Semiconductor devices are used in a variety of electronic applications, such as personal computers, cell phones, digital cameras, and other electronic equipment, as examples. Semiconductor devices are typically fabricated by sequentially depositing insulating or dielectric layers, conductive layers, and semiconductive layers of material over a semiconductor substrate, and patterning the various material layers using lithography to form circuit components and elements thereon.\nThe semiconductor industry continues to improve the integration density of various electronic components (e.g., transistors, diodes, resistors, capacitors, etc.) by continual reductions in minimum feature size, which allow more components to be integrated into a given area. These smaller electronic components also require smaller packages that utilize less area than packages of the past, in some applications.\nOne smaller type of packaging for semiconductors is a flip chip chip-scale package (FcCSP), in which a semiconductor die is placed upside-down on a substrate and bonded to the substrate using bumps. The substrate has wiring routed to connect the bumps on the die to contact pads on the substrate that have a larger footprint. An array of solder balls is formed on the opposite side of the substrate and is used to electrically connect the packaged die to an end application.\nHowever, some FcCSP packages tend to exhibit bending, where warping of the substrate occurs during processing, such as during temperature stress. The bending can cause reliability issues, such as bond breakage of the bumps, delamination of an underfill, and delamination of a passivation layer on the die."}
-{"text": "Polar organisms should overcome the problems of decreased enzyme activity, decreased membrane fluidity, inactivation and improper folding of proteins, formation of intracellular ice crystals, etc. to survive in low-temperature, polar environments. Among others, the formation of ice crystals causes physical damages and dehydration of tissues due to the growth of ice crystals, thus causing serious damage to polar organisms. Polar organisms produce various antifreeze proteins (hereinafter referred to as \u201cAFPs\u201d) to survive at low temperatures. AFPs inhibit the growth of ice crystals in vivo and the recrystallization of ice to protect polar organisms from sub-zero temperatures to survive (Davies, P. L. and Sykes, B. D., Curr. Opin. Struct. Biol. 7, 1997, 828-834; Davies, P. L. et al., Philos Trans R Soc Lond B Biol Sci. 357, 2002, 927-935; D'Amico, S. et al., EMBO Rep. 7, 2006, 385-389).\nAFPs are proteins that generally have a flat ice-binding surface and bind to specific surfaces of ice crystals, thus inhibiting the growth of ice crystals and the recrystallization of ice. AFPs create a difference between the melting point and freezing point. This is called thermal hysteresis (TH), which can be measured using a nanoliter osmometer and used as an indicator of AFP activity. Moreover, AFPs do not lower the freezing point in proportion to the concentration, unlike typical antifreeze used in vehicles. That is, AFPs can effectively lower the freezing point even at very low concentrations by direct interaction with ice, thus minimizing damage due to osmotic pressure generated in vivo during freezing (Jia, Z. and Davies P. L., Trends Biochem. Sci. 27, 2002, 101-106).\nThe unique features of AFPs that prevent the growth of ice crystals and inhibit the recrystallization of ice have been used in various commercial fields. For example, in the agricultural field, AFP expression in plants has been attempted for the purpose of preventing cold-weather damage to plants. Moreover, in the field of fisheries, there has been an attempt to produce a transgenic fish by expressing AFPs in commercially available fish such as Atlantic salmon (Salmo salar) or goldfish (Carassius auratus) so as to enable farming in cold areas. Furthermore, in the medical field, research on the use of AFPs in cryosurgery and as an additive in cryopreservation of blood, stem cells, umbilical cord blood, organs, and germ cells has continued to progress. In addition, in the food field, AFPs are also used in product production for frozen storage of smoother ice scream. In the field of cosmetics, functional cosmetics containing AFPs for preventing frostbite have already been sold. Although AFPs are widely used in various commercial fields as mentioned above, there are still limitations in mass production of recombinant AFPs due to low-level expression of AFPs and folding problems. This is mainly because most AFPs have disulfide bonds and are stabilized by disulfide bonds, which thus makes it difficult to express recombinant proteins and yields improper folding of expressed proteins.\nSince AFPs were first discovered in fish living in cold water, various types of new AFPs have been discovered in insects, plants, fungi, microorganisms, etc. New AY30 AFP derived from arctic yeast, Leucosporidium sp., has recently been recovered. The AY30 AFP has no cysteine amino acid residues, and thus during production of recombinant proteins, the level of protein expression is high, and the folding problem due to improperly formed disulfide bonds does not occur, As a result, the AY30 AFP is suitable for mass production of recombinant AFPs.\nTherefore, the present inventors have synthesized a recombinant polynucleotide by modifying an AFP gene to be expressed using codon optimization for a yeast expression system and inserted the recombinant polynucleotide into a yeast-derived expression vector so as to mass-produce an antifreeze protein (AFP) derived from arctic yeast by overexpressing AFP in the form of activated protein. As a result, the present inventors have obtained a large amount of AFP and found that the AFP is glycosylated, thus completing the present invention. All references cited in this specification are hereby incorporated by reference in their entirety."}
-{"text": "Plasma arc torches are widely used in the cutting, and marking of materials. A plasma torch generally includes an electrode and a nozzle having a central exit orifice mounted within a torch body, electrical connections, passages for cooling, and passages for arc control fluids (e.g., plasma gas). Optionally, a swirl ring is employed to control fluid flow patterns in the plasma chamber formed between the electrode and nozzle. In some torches, a retaining cap can be used to maintain the nozzle and/or swirl ring in the plasma arc torch. The torch produces a plasma arc, a constricted ionized jet of a gas with high temperature and high momentum. Gases used in the torch can be non-reactive (e.g., argon or nitrogen) or reactive (e.g., oxygen or air). In operation, a pilot arc is first generated between the electrode (cathode) and the nozzle (anode). Generation of the pilot arc can be by means of a high frequency, high voltage signal coupled to a DC power supply and the torch or by means of any of a variety of contact starting methods.\nOne category of hand held plasma arc torch systems include a manual gas control knob on the control panel of the power supply or power supply housing. Before cutting a workpiece, an operator is required to manually adjust the gas pressure or gas flow rate based on the process parameters set forth in a cut chart. The operator manually adjusts the gas pressure or flow rate for each type of cut and therefore, constantly refers to the cut chart for the appropriate gas pressure or flow rate. Moreover, if the operator inadvertently inputs an incorrect gas pressure or flow rate, the plasma arc torch can operate incorrectly or can operate inefficiently.\nAnother category of hand held systems eliminate the gas control by automatically setting the gas pressure based on the user selected current level and mode (i.e., gouging or cutting). This category of hand held plasma arc torches does not provide the operator with any flexibility in setting the gas pressure beyond the preset automated systems. Therefore, if the operator determines that the gas pressure or flow rate should be changed due to a changed operating parameter or to optimize the plasma arc torch, the operator does not have the flexibility to make these operational and/or optimizing adjustments."}
-{"text": "There is a constant rise in demand for artificial meniscal grafts mimicking native articular tissue to be used for surgical treatment of meniscal lesions. In Europe alone over 400,000 surgical cases involving the meniscus are being performed annually, and over 1 million similar cases are treated in the United States. By far meniscectomy is known to be the most common surgical procedure performed in the orthopedic field today. The current therapeutic strategy for this type of meniscus tears is either partial or subtotal meniscectomy, with only a small percentage being successfully repaired but finally leading to osteoarthritis of the knee with time (Fairbank, 1948; Englund et al., 2003).\nA functional intact meniscus is of paramount importance for homeostasis of the knee joint. It helps perform complex knee joint biomechanics, in load bearing, load transmission, shock absorption, joint stability and joint lubrication. However, due to lack of vasculature, human meniscus has a poor healing potential. Blood vessels are reported to be present only in the outer 10-30% of the meniscal body and can be sutured successfully with a high success rate (Englund et al., 2003; Buma et al., 2004). In contrast, majority of these meniscal tears are situated in the inner avascular zone lacking spontaneous healing process and hence be resected (Kohn et al., 1999). Removal and/or damage of all this important anatomical structure eventually leads to degenerative changes of the articular cartilage, osteoarthritis and subsequent clinical symptoms due to increased peak stresses (Fairbank, 1948; Cole et al., 2003; Chatain et al., 2003; Englund et al., 2003). It has been estimated that cartilage volume loss after meniscectomy is at 4% per year and is known to be more pronounced in the lateral compartment as compared to medial compartment (Verdonk and Kohn, 1999).\nTo this problem, meniscus allo/autograft transplantation represents a potential tissue engineering solution for the symptomatic, meniscus deficient patient to substitute for lost meniscal tissue to prevent cartilage degeneration, relieve pain and to improve function. The strategies included delivery of potent cells to the defect site for repair including chondrocytes, fibrochondrocytes and stem cells (Peretti et al., 2004; Izuta et al., 2005; Port et al., 1996). The other strategy being direct replacement of defective tissue in part or as a whole has also been carried out using both natural and synthetic scaffolds, including collagen-based grafts, subintestinal submucusa, cell free hydrogels, degradable porous foams, macro- and microporous polymeric meshes to improve immediate or long term outcomes (Buma et al., 2004; Stone et al., 1992; Cook et al., 2006 a; Setton et al., 1999; Sweigart et al., 2001; Kobayashi et al., 2005; Kelly et al., 2007; Van Tienen et al., 2002; Heijkants et al., 2004; Cook et al., 2006 (a, b). In the past, a variety of these materials have already been reported for cartilage tissue engineering including, poly-glycolic acid (PGA), poly-L-lactic acid (PLA), copolymer poly-lactic-co-glycolic acid (PLGA) and alginate (Grande et al., 1997; Freed et al., 1993 a,b; Paige et al., 1996; Marijnissen et al., 2002; Ma et al., 2003). However, these materials have intrinsic limitations, including inflammation in vivo in the case of the polyesters and rapid degradation and high swelling in the case of collagen, which can limit their use (Cancedda et al., 2003; Athanasiou et al., 1996; Wakitani et al., 1994; Meinel et al., 2004 a,b). In terms of meniscus shape, a PGA spun matrix was used in a rabbit model but failed to recapitulate the complex internal meniscus architecture (Kang et al., 2006). Additional efforts have focused on mimicking the native mesh-like meniscus architecture using cell alignment on biodegradable electrospun fibers for enhanced biomechanics (Baker and Mauck; 2007; Baker et al., 2009). Many of the above studies employed in vivo animal models to show chondroprotection by the implant, but with a low success rate due to failure to mimic the complex internal architecture and biomechanics of the native meniscus.\nIn order to develop a functional tissue engineered meniscus, mimicking its complex internal architecture is most important. In this regard, none of the approaches previously reported have successfully recapitulated the complex native meniscal multiporous and aligned structure as a single meniscus wedge shaped unit to completely and/or partially eliminate cartilage degeneration. Thus, in order to mimic the meniscus in a tissue engineered approach, understanding its structural and functional components is important. Menisci are wedge-shaped semi-lunar discs present in duplicate in each knee joint which are attached to the transverse ligaments, the joint capsule, the medial collateral ligament (medially) and the menisco-femoral ligament (laterally) (McDevitt and Webber, 1990; Sweigart and Athanasiou, 2001). An extensive scanning electron micrograph study of the human meniscus by Peterson and Tillmann showed 3 distinct zones comprising of outer finer meshwork, middle broader mesh like fibrous structure and bottom most aligned collagen bundles in laminar orientation (Petersen and Tillmann, 1998). This particular aligned laminar orientation of fibers along with mesh structure within was reported to contribute maximally for its high intrinsic tensile and compressive properties of native meniscus (Sweigart and Athanasiou, 2001; Tissakht and Ahmed, 1995; Petersen and Tillmann, 1998). As a fibrocartilaginous structure, the meniscus has characteristic of both fibrous (outer region) and cartilaginous (inner region) properties (O'Connor, 1976; Petersen and Tillmann, 1998). Knee meniscal fibrocartilaginous tissue contains mainly water (72%), collagens (22%) and glycosaminoglycans (0.8%) (Proctor et al, 1989; Herwig et al, 1984). Of the total collagen content, Type I collagen accounts for over 90%. The remaining 10% are meniscal collagens Type II, III and V collagen (Eyre and Wu, 1983; McDevitt and Webber, 1990). It has been shown that peripheral two-thirds of the meniscus solely consist of type I collagen, whereas type II collagen comprises a large portion of the fibrillar collagen on the inner side (Cheung, 1987). Proteoglycans make for 2-3% of the dry weight and are mainly concentrated in the inner cartilaginous region of the meniscus (McDevitt and Webber, 1990; Buma et al., 2004). Also, the cellular component of the meniscus further reflects its fibrocartilaginous nature, the main cell type being meniscus fibrochondrocytes (McDevitt and Webber, 1990). Regarding cell types, at least two cell populations are present within the human meniscus (Ghadially et al., 1983). The fibrochondrocytes being the main cell type are reported within the inner and middle part of the meniscus having a rounded or oval shaped cell structure surrounded by an abundant ECM deposition (McDevitt and Webber, 1990; Ghadially et al., 1983). The outer one-third meniscus is reported to be populated mainly by spindle shaped fibroblast like cells with a dense connective tissue (Ghadially et al., 1983).\nOver the years, newer improvised methods such as meniscus allograft or autograft transplantation have been constantly searched for substituting the resected meniscus in case of either total or partial meniscectomy. However, none to date have generally been able to recapitulate and recreate the native meniscal multiporous and aligned structure as a single meniscus wedge shaped unit to completely and/or partially eliminate cartilage regeneration. As such, there is still a strong need to develop a scaffold that can mimic heterogeneous architecture and functions of native meniscal tissue."}
-{"text": "1. Technical Field\nThe relates to a phase-change random access memory (PCRAM) device, and more particularly, to a PCRAM device and a method of manufacturing the same.\n2. Related Art\nWith demands on lower power consumption, next-generation memory devices having nonvolatile and non-refresh properties have been studied. A PCRAM device of the next-generation memory devices includes a switching element connected at intersections of word lines and bit lines, which are arranged to cross each other, a lower electrode electrically connected to the switching element, a phase-change layer formed on the lower electrode, and an upper electrode formed on the phase-change layer.\nIn a conventional PCRAM device, when a write current flows through the switching element and the lower electrode, Joule heat is generated at an interface between the phase-change layer and the lower electrode. The phase-change layer is phase-changed into an amorphous state or a crystalline state by the generated joule heat. Therefore, the conventional PCRAM device stores data using a difference between resistances in the amorphous state and the crystalline state of the phase-change layer.\nHowever, in the conventional PCRAM device, the Joule heat generated when the write current flows affects a phase-change layer of adjacent cell.\nThe effect on adjacent cells is generally referred to as thermal disturbance. In recent years, the thermal disturbance has an increased effect on adjacent cells when a semiconductor memory device is highly integrated.\nFIGS. 1A and 1B are views illustrating thermal disturbance of a conventional PCRAM device.\nAs shown in FIGS. 1A and 1B, the conventional PCRAM device includes a lower electrode 10 formed on a switching element (not shown), a phase-change layer 20 formed on the lower electrode 10, and an upper electrode 30 formed on the phase-change layer 20. The reference numeral 40 denotes an insulating layer.\nAs shown in FIG. 1A, if a cell A is written when cells B are written with data \u201c1\u201d, which is a high resistance state, Joule heat is generated at an interface between the lower electrode 10 and the phase-change layer 20 of the cell A (see FIG. 1B), and thus, phase-change material patterns of amorphous states in the cells B are crystallized. Therefore, resistances of the cells B are reduced.\nThe thermal disturbance generated in the conventional PCRAM device may cause a malfunction, and thus reliability of the conventional PCRAM device is degraded."}
-{"text": "This invention relates to the simultaneous measurement of the concentration of a selected ion species in a solution and the pH of the solution. The invention particularly, though not exclusively, relates to photographic solutions, and particularly, though not exclusively, where the selected ion species is silver. In general, however, the invention relates to the simultaneous potentiometric measurement of the concentration of any ion species in a solution and measurement of the pH of the solution using an ISFET (Ion Selective Field Effect Transistor).\nFor the present purpose the tern xe2x80x9csolutionxe2x80x9d is to be understood as also including an emulsion, for example a mixture of a silver compound suspended in gelatin, or a dispersion. The invention will be particularly described, by way of example only, with reference to photographic solutions.\nIt is known simultaneously to measure silver ion concentration in, and the pH of, an aqueous solution. In one arrangement, a single reference electrode is connected into a first potentiometer circuit with a conventional glass pH electrode, and is connected into a second potentiometer circuit with a conventional silver electrode, all three electrodes being immersed in the solution. In another arrangement, an ISFET is used instead of the glass pH electrode. This necessitates the use of a separate reference electrode for each measuring circuit in order to provide electrical isolation between the circuits since the ISFET is a current carrying device whose presence would otherwise interfere with the voltage measurement of the silver electrode.\nA glass pH electrode has the disadvantage that it can be damaged under conditions of high temperature and high pH, so that its readings become unreliable or inconsistent. An ISFET overcomes this disadvantage. However, the conventional arrangement including an ISFET described above is complicated by the requirement of the additional reference electrode, especially when applied in a large scale production vessel, as used in the preparation of photographic emulsions for example, where the electrodes are configured in a unitary probe. This can lead to difficulties for maintenance and for calibration. Furthermore, existing probe structures would require extensive modification to accommodate the additional reference electrode, which would be expensive.\nIt will be appreciated that if, on the other hand, measurement of ion concentration and pH were not required simultaneously, then the measurements would not interfere with each other and a single reference electrode could be used successively in combination with an ion concentration electrode and an ISFET.\nIn accordance with one aspect of the present invention, there is provided apparatus for simultaneously measuring the concentration of a selected ion species in a solution and the pH of the solution, comprising: a first electrical circuit that is arranged to receive signals from both a reference electrode and an ion selective electrode immersed in the solution and to derive therefrom an output signal representative of the concentration of the selected ion in the solution; a second electrical circuit that is arranged to receive signals from both said reference electrode and an ISFET immersed in the solution and to derive therefrom an output signal representative of the pH of the solution; wherein any d.c. input signal to said first electrical circuit from the reference electrode is substantially electrically isolated from the input of the second circuit; wherein a signal representative of the voltage, usually earth potential, of the solution is supplied (a) directly to the first circuit so as to establish a reference, usually earth, potential for the first circuit, and (b) to the second circuit through a.c. coupling means so as to establish a corresponding virtual reference, usually earth, potential for the second circuit; and wherein the first and second electrical circuits are arranged to be provided with electrical power from supplies that are electrically isolated from each other.\nThe apparatus may comprise means for displaying a representation of said ion concentration and pH output signals, wherein said second electrical circuit includes an isolation amplifier, and wherein said display means is arranged to receive said pH output signal of the second circuit through the isolation amplifier. Preferably, the apparatus includes a further isolation amplifier through which the ion concentration output signal of the first circuit is supplied to the display means. Advantageously, the apparatus comprises a low pass filter, wherein said pH output signal from the second electrical circuit is arranged to be passed to the display means through the low pass filter.\nPreferably, the apparatus comprises a high value resistor, for example of about 1 Mxcexa9 or greater, that is arranged to effect said electrical isolation of d.c. input signals to said first and second electrical circuits. Also said a.c. coupling means may comprise a high value capacitor, for example of about 1 xcexcF or greater.\nIn accordance with another aspect of the present invention, there is provided a method of simultaneously measuring the concentration of a selected ion species in a solution and the pH of the solution, comprising the steps of: measuring in a first electrical circuit the potential difference between an ion selective electrode and a reference electrode both immersed in the solution, and deriving therefrom the concentration of the ions in the solution; measuring in a second electrical circuit the current flowing between an ISFET and the reference electrode both immersed in the solution, and deriving therefrom the pH of the solution; connecting the reference electrode to the first and second electrical circuits such that any d.c. signal from the reference electrode is electrically isolated from the second circuit; making an electrical connection between the solution and the first circuit so as to provide the solution potential as a reference, preferably earth, potential therefore, and making an electrical connection between the solution and the second circuit through a.c. coupling means so as to provide a corresponding virtual reference, preferably earth, potential therefore; and supplying the first and second circuits with electrical power from sources that are electrically isolated from each other.\nThe method of the invention is advantageously carried out using the apparatus of the invention.\nDetails of electrodes suitable for use in the present invention as ion selective and reference electrodes, and of ISFETs, can be found in the book xe2x80x9cpH Measurementxe2x80x9d by Helmuth Galster (VCH,1991).\nThe electrical isolation of the two circuits provided in the present invention allows an ISFET to be used in the pH measuring circuit, whilst needing only a single, common, reference electrode. The disadvantages of the known arrangements for simultaneous ion concentration and pH measurement are thus overcome in a particularly convenient manner.\nThe isolation is provided at several stages. Initially this is done by arranging that the signal from the reference electrode is used in the ion concentration circuit as a potentiometric measurement, and is supplied to the pH measuring circuit only as an a.c. input, i.e. after having any d.c. component isolated therefrom. An actual reference potential, the potential, usually earth, of the solution, is applied to the first circuit, and a virtual reference potential derived therefrom is applied to the second circuit. The two circuits have separate isolated power supplies. Furthermore, when the resulting ion concentration and pH signals are supplied to a display means, such as a multi-channel voltmeter, this is done through respective isolation amplifiers, which are preferably supplied from a third, isolated power supply.\nThe ability to use a single reference electrode means that a single, unitary measurement probe can be constructed, in which the ISFET can be installed relatively easily along with the ion selective and reference electrodes. Furthermore, the measuring apparatus can be calibrated more easily than is the case with the known arrangement using two reference electrodes.\nAlthough described with reference to a single ion selective electrode and a single ISFET, it is envisaged that the present invention may comprise two or more ion concentration electrodes and/or two or more ISFETs, each type of electrode being connected into the respective first or second electrical circuit.\nApparatus for, and a method of, simultaneously measuring the concentration of a selected ion species in a solution and the pH of the solution, will now be described, by way of example, with reference to the accompanying schematic circuit diagram."}
-{"text": "1. Field of the Invention\nThe invention relates to a diffractive optical element having a multiplicity of binary blazed diffraction structures. The diffractive optical element is particularly intended for use in microlithographic projection exposure apparatus.\n2. Description of the Prior Art\nConventional blazed gratings have diffraction structures of triangular, in particular sawtoothed cross section which extend mutually parallel with a spacing equal to the grating constant g. One edge of the diffraction structures, the blaze edge, has an inclination with respect to the base surface of the grating such that the reflection or refraction law is satisfied for one diffraction order of the incident light, and the majority of the intensity of the diffracted light is therefore contained in the order favoured by the blaze edge. The traditional method of producing such blazed gratings consisted in scratching the diffraction structures in a master grating with the aid of diamonds and making corresponding copies of this master grating. This mechanical method is highly elaborate, on the one hand, and on the other hand it encounters limitations with very short wavelengths of the light for which the grating is intended to be used, since the structures to be produced are too small.\nEfforts have therefore been made to employ the process technology used for the production of semiconductor components, in which a substrate is coated with photoresist, exposed, subsequently developed and etched, in order to produce the diffraction structures of blazed gratings. The approach firstly involved using successions of such process cycles to achieve diffraction structures which are supposed to approximate the blaze edge by a stepped edge. If four such steps are used, for example, then diffraction efficiencies of more than 80% can be achieved in the first order. With a further process cycle, eight stages are obtained by which a first-order diffraction efficiency of about 95% can be achieved. In general, 2n steps can be produced by using n process cycles. With increasing n, the stepped profile of the edge becomes closer and closer to the sawtooth profile of ideal blazed gratings in conventional, mechanically produced gratings, the diffraction efficiency of which is 100% in the first order according to scalar theory. The production of such a grating, however, is cost-intensive and error-prone because it is necessary to carry out the process cycle repeatedly.\nAttempts have also been undertaken to simulate the blaze profile of the diffraction structures by using binary structures whose dimensions are smaller than the wavelength of the electromagnetic radiation for which the grating was defined. These attempts are based on the fact that light is no longer diffracted at the small substructures, but can only be scattered. This leaves only the zeroth diffraction order which picks up the effect of the substructures merely in the form of a local effective refractive index in phase gratings, or merely in the form of a local shade of grey in amplitude gratings.\nA first example of such a binary blazed grating is described in the article by Joseph N. Mait et al. \u201cDiffractive lens fabricated with binary features less than 60 nm\u201d, Optics Letters, 15 Mar. 2000, pages 381 et seqq. The substructure used here is a multiplicity of lines, all of which extend parallel to the diffraction structure and whose spacing is less than the effective wavelength.\nThe article by Philippe Lalanne et al. \u201cDesign and fabrication of blazed binary diffractive elements with sampling periods smaller than the structural cut off\u201d, J. Opt. Soc. Am. A, May 1999, pages 1143 et seqq. describes blazed diffractive elements of the type mentioned in the introduction, in which the diffraction structures are resolved into individual substructures consisting of rectangular or square pillars. Different \u201cfill factors\u201d can be achieved by varying the pillar width for a predetermined pillar spacing, and this corresponds to a local variation of the effective refractive index. As an alternative, the pillars may also be arranged at different spacings with a constant width.\nA common feature of all these attempts to produce binary blazed diffractive optical elements is that the substructures are minutely configured and have a very high aspect ratio (structure height to structure width). They are therefore technologically highly elaborate and expensive to produce, and cannot be made with sufficient accuracy."}
-{"text": "This invention relates generally to turbine engine stator assemblies, and more particularly, to apparatus and method for controlling operating clearance between a stationary shroud surface in a turbine engine stator assembly and a rotating surface of juxtaposed blading members.\nForms of an axial flow turbine engine, typically a gas turbine engine, include rotating assemblies radially within stationary assemblies that assist in defining a flowpath of the engine. Examples include a rotary compressor assembly that compresses incoming air, and a rotary turbine assembly that extracts power from products of engine fuel combustion. Such assemblies comprise stages of rotating blades within a surrounding stator assembly that includes a shroud surface spaced apart from cooperating surfaces of the rotating blades. Efficiency of a turbine engine depends, at least in part, on the clearance or gap between the juxtaposed shroud surface and the rotating blades. If the clearance is excessive, undesirable leakage of engine flowpath fluid will occur between such gap resulting in reduced engine efficiency. If the clearance is too small, interference can occur between the rotating and stationary members of such assemblies, resulting in damage to one or more of such cooperating surfaces.\nComplicating clearance problems in such apparatus is the well known fact that clearance between such turbine engine assemblies changes with engine operating conditions such as acceleration, deceleration, or other changing thermal or centrifugal force conditions experienced by the cooperating members during engine operation. Clearance control mechanisms for such assemblies, sometimes referred to as active clearance control systems, have included mechanical systems or systems based on thermal expansion and contraction characteristics of materials for the purpose of maintaining selected clearance conditions during engine operation. Such systems generally require use of substantial amounts of air for heating or cooling at the expense of such air otherwise being used in the engine operating cycle. Provision of an improved means for active clearance control that reduces the need for engine flowpath fluid for such heating or cooling could enhance engine efficiency."}
-{"text": "1. Field of the Invention\nThis invention relates generally to a vehicle rollover avoidance system and, more particularly, to a vehicle rollover avoidance system that employs a roll control factor and a yaw rate stability control factor to control semi-active suspension dampers to mitigate the risk of vehicle rollover.\n2. Discussion of the Related Art\nIt is known in the art to mitigate a potential vehicle rollover using differential braking control, rear-wheel steering control, front-wheel steering control, or any combination thereof. A vehicle rollover avoidance system may receive vehicle dynamics information from various sensors, such as yaw rate sensors, lateral acceleration sensors and roll rate sensors, to determine the proper amount of action to be taken to detect a potential vehicle rollover. A balance typically needs to be provided between estimating the vehicle roll motion and the vehicle yaw motion to provide the optimal vehicle response. Thus, it is usually necessary to detect certain vehicle conditions to provide the roll detection. To precisely identify vehicle roll stability conditions, it may be advantageous to know the vehicle's roll rate and roll angle because they are the most important states in vehicle roll dynamics.\nUnder normal driving conditions, drivers can direct the vehicle to the desired heading through the control of the steering wheel. When the vehicle is turning, there are actually three motions taking place with the vehicle. Particularly, a turning motion, or yaw, is occurring, as the vehicle body is turning around an imaginary access vertical to the ground through the so-called yaw-center of the vehicle. Also, there is subtle vehicle sliding laterally, sometimes in the direction of the turn and sometimes away from the turn, depending mainly on the vehicle speed. Further, a tilting motion or roll motion occurs as the vehicle's body is turning around an imaginary axis parallel to the ground through the so-called roll-axis of the vehicle.\nUnder normal vehicle maneuvering conditions, the tire/road contact surfaces can generate sufficient forces to sustain the desired vehicle motions, and drivers are accustomed with these motions as they occur. However, when the vehicle maneuver starts approaching limit-handling conditions, the tire/road contact surfaces can no longer sustain the desired yaw motion and side-slip motion, and the vehicle body will exhibit an increased roll motion. As a result, a discrepancy will build up between the vehicle's yaw rate and its desired yaw rate, and between the vehicle's side-slip velocity and its desired side-slip velocity. Further, if the roll motion becomes too large, the vehicle may roll over."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to the design of integrated circuits and more particularly to sense amplifiers.\n2. Description of the Background Art\nMany systems on an integrated circuit are designed to respond differently depending upon whether their input voltages are considered HIGH or LOW. Sometimes, an input voltage must be modified to conform to the HIGH or LOW state (e.g., during the period when the input voltage transitions between states). Sense amplifiers are circuits that detect a small voltage differential and increase or decrease the voltage to a level required by the system. An example of a system that utilizes sense amplifiers is a computer memory circuit. Information stored in the memory cells of a memory chip using sense amplifiers can be retrieved much faster than from a memory chip without sense amplifiers.\nAs shown in FIG. 1, a common static random access memory (SRAM) configuration generally designated 100 includes an array 105 of memory cells 110. Each memory cell 110 is connected to a word line 115, a bit line B 120, and a complement of the bit line, B 145. The memory cells 110 connected to each of the word lines 115 define a memory cell array row 125, and the memory cells connected to each of the bit line B 120 and a corresponding complement of the bit line B 145 define a memory cell array column 130. Each memory cell 110 stores information in the form of a voltage charge representing a logic value of LOW or HIGH. A voltage level equal to V.sub.DD represents the logic value of HIGH and V.sub.SS represents the logic value of LOW.\nBit lines B 120 and B 145 are connected to an equalization and precharge circuit 150. The precharge component of the equalization and precharge circuit 150 initially charges bit lines B 120 and B 145 to the voltage level of V.sub.DD. The equalization component of the equalization and precharge circuit 150 ensures that the voltages on bit lines B 120, .nu..sub.B, and B 145, .nu..sub.B, are initially at the same level.\nThe word lines 115 are connected to a row decoder 155. When a memory cell 110' is accessed, the row decoder 155 selects and changes the voltage of a word line 115' corresponding to memory cell 110'. A changed voltage signal (e.g., LOW to HIGH) from the word line 115' allows memory cell 110' to communicate with bits lines B 120' and B 145'. If memory cell 110' stores a logic value of HIGH, then .nu..sub.B will remain at HIGH and .nu..sub.B will decrease to LOW. If memory cell 110' stores a logic value of LOW, then .nu..sub.B will decrease to LOW and .nu..sub.B will remain at HIGH.\nBit lines B 120 and B 145 are connected to a sense amplifier 160 which detects and amplifies the difference in voltage between .nu..sub.B and .nu..sub.B. Depending on the difference between .nu..sub.B and .nu..sub.B, the sense amplifier 160 will output either V.sub.DD or V.sub.SS.\nConnected to the sense amplifier 160 is a column decoder 165. The column decoder 165, like the row decoder 155, includes a combination of logic circuits that select a logic signal from either one or a set of the memory cell array columns 130 for final output from SRAM 100.\nThe prior art described above suffers from a number of limitations. To store more information on a single memory chip, smaller memory cells are used. Smaller memory cells, however, use smaller transistors, which have less driving capability, resulting in a longer time for .nu..sub.B and .nu..sub.B to reach distinct HIGH or LOW voltage levels. To reduce the time required to read a memory cell, sense amplifiers are used to quickly detect the small voltage difference between .nu..sub.B and .nu..sub.B without having to wait for .nu..sub.B and .nu..sub.B to reach definite HIGH or LOW voltage levels. However, when .nu..sub.B and .nu..sub.B reach definite HIGH or LOW voltage levels before the operation of the sense amplifier, the operation of the sense amplifier is not required and consumes unnecessary power.\nWhat is needed is a sense amplifier design that overcomes the shortfalls of the sense amplifier designs known in the art."}
-{"text": "Certain embodiments of the present invention are directed to integrated circuits. More particularly, some embodiments of the invention provide a system and method for stage-based control related to TRIAC dimmer. Merely by way of example, some embodiments of the invention have been applied to driving one or more light emitting diodes (LEDs). But it would be recognized that the invention has a much broader range of applicability.\nA conventional lighting system may include or may not include a TRIAC dimmer that is a dimmer including a Triode for Alternating Current (TRIAC). For example, the TRIAC dimmer is either a leading-edge TRIAC dimmer or a trailing-edge TRIAC dimmer. Often, the leading-edge TRIAC dimmer and the trailing-edge TRIAC dimmer are configured to receive an alternating-current (AC) input voltage, process the AC input voltage by clipping part of the waveform of the AC input voltage, and generate a voltage that is then received by a rectifier (e.g., a full wave rectifying bridge) in order to generate a rectified output voltage.\nFIG. 1 shows certain conventional timing diagrams for a leading-edge TRIAC dimmer and a trailing-edge TRIAC dimmer. The waveforms 110, 120, and 130 are merely examples. Each of the waveforms 110, 120, and 130 represents a rectified output voltage as a function of time that is generated by a rectifier. For the waveform 110, the rectifier receives an AC input voltage without any processing by a TRIAC dimmer. For the waveform 120, an AC input voltage is received by a leading-edge TRIAC dimmer, and the voltage generated by the leading-edge TRIAC dimmer is received by the rectifier, which then generates the rectified output voltage. For the waveform 130, an AC input voltage is received by a trailing-edge TRIAC dimmer, and the voltage generated by the trailing-edge TRIAC dimmer is received by the rectifier, which then generates the rectified output voltage.\nAs shown by the waveform 110, each cycle of the rectified output voltage has, for example, a phase angel (e.g., \u03d5) that changes from 0\u00b0 to 180\u00b0 and then from 180\u00b0 to 360\u00b0. As shown by the waveform 120, the leading-edge TRIAC dimmer usually processes the AC input voltage by clipping part of the waveform that corresponds to the phase angel starting at 0\u00b0 or starting at 180\u00b0. As shown by the waveform 130, the trailing-edge TRIAC dimmer often processes the AC input voltage by clipping part of the waveform that corresponds to the phase angel ending at 180\u00b0 or ending at 360\u00b0.\nVarious conventional technologies have been used to detect whether or not a TRIAC dimmer has been included in a lighting system, and if a TRIAC dimmer is detected to be included in the lighting system, whether the TRIAC dimmer is a leading-edge TRIAC dimmer or a trailing-edge TRIAC dimmer. In one conventional technology, a rectified output voltage generated by a rectifier is compared with a threshold voltage Vth_on in order to determine a turn-on time period Ton. If the turn-on time period Ton is approximately equal to the duration of a half cycle of the AC input voltage, no TRIAC dimmer is determined to be included in the lighting system; if the turn-on time period Ton is not approximately equal to but is smaller than the duration of a half cycle of the AC input voltage, a TRIAC dimmer is determined to be included in the lighting system. If a TRIAC dimmer is determined to be included in the lighting system, a turn-on voltage slope Von_slope is compared with the threshold voltage slope Vth_slope. If the turn-on voltage slope Von_slope is larger than the threshold voltage slope Vth_slope, the TRIAC dimmer is determined to be a leading-edge TRIAC dimmer; if the turn-on voltage slope Von_slope is smaller than the threshold voltage slope Vth_slope, the TRIAC dimmer is determined to be a trailing-edge TRIAC dimmer.\nIf a conventional lighting system includes a TRIAC dimmer and light emitting diodes (LEDs), the light emitting diodes may flicker if the current that flows through the TRIAC dimmer falls below a holding current that is, for example, required by the TRIAC dimmer. As an example, if the current that flows through the TRIAC dimmer falls below the holding current, the TRIAC dimmer may turn on and off repeatedly, thus causing the LEDs to flicker. As another example, the various TRIAC dimmers made by different manufacturers have different holding currents ranging from 5 mA to 50 mA.\nThe light emitting diodes (LEDs) are gradually replacing incandescent lamps and becoming major lighting sources. The LEDs can provide high energy efficiency and long lifetime. The dimming control of LEDs, however, faces significant challenges because of insufficient dimmer compatibility. For certain historical reasons, the TRIAC dimmers often are designed primarily suitable for incandescent lamps, which usually include resistive loads with low lighting efficiency. Such low lighting efficiency of the resistive loads often helps to satisfy the holding-current requirements of TRIAC dimmers. Hence the TRIAC dimmers may work well with the incandescent lamps. In contrast, for highly efficient LEDs, the holding-current requirements of TRIAC dimmers usually are difficult to meet. The LEDs often need less amount of input power than the incandescent lamps for the same level of illumination.\nIn order to meet the holding-current requirements of the TRIAC dimmers, some conventional techniques use a bleeder for a lighting system. FIG. 2 is a simplified diagram of a conventional lighting system that includes a bleeder. As shown, the conventional lighting system 200 includes a TRIAC dimmer 210, a rectifier 220, a bleeder 224, a diode 226, capacitors 230, 232, 234, 236 and 238, a pulse-width-modulation (PWM) controller 240, a winding 260, a transistor 262, resistors 270, 272, 274, 276, 278 and 279, and one or more LEDs 250. The PWM controller 240 includes controller terminals 242, 244, 246, 248, 252, 254, 256 and 258. For example, the PWM controller 240 is a chip, and each of the controller terminals 242, 244, 246, 248, 252, 254, 256 and 258 is a pin. In yet another example, the winding 260 includes winding terminals 263 and 265.\nThe TRIAC dimmer 210 receives an AC input voltage 214 (e.g., VAC) and generates a voltage 212. The voltage 212 is received by the rectifier 220 (e.g., a full wave rectifying bridge), which then generates a rectified output voltage 222. The rectified output voltage 222 is larger than or equal to zero. The resistor 279 includes resistor terminals 235 and 239, and the capacitor 236 includes capacitor terminals 281 and 283. The resistor terminal 235 receives the rectified output voltage 222. The resistor terminal 239 is connected to the capacitor terminal 281, the controller terminal 252, and a gate terminal of the transistor 262. The gate terminal of the transistor 262 receives a gate voltage 237 from the resistor terminal 239, the capacitor terminal 281, and the controller terminal 252. The capacitor terminal 283 receives a ground voltage.\nAs shown in FIG. 2, the rectified output voltage 222 is used to charge the capacitor 236 through the resistor 279 to raise the gate voltage 237. In response, if the result of the gate voltage 237 minus a source voltage at a source terminal of the transistor 262 reaches or exceeds a transistor threshold voltage, the transistor 262 is turned on. When the transistor 262 is turned on, through the transistor 262 and the controller terminal 254, a current flows into the PWM controller 240 and uses an internal path to charge the capacitor 232. In response, the capacitor 232 generates a capacitor voltage 233, which is received by the controller terminal 244. If the capacitor voltage 233 reaches or exceeds an undervoltage-lockout threshold of the PWM controller 240, the PWM controller 240 starts up.\nAfter the PWM controller 240 has started up, a pulse-width-modulation (PWM) signal 255 is generated. The PWM signal 255 has a signal frequency and a duty cycle. The PWM signal 255 is received by the source terminal of the transistor 262 through the terminal 254. The transistor 262 is turned on and off, in order to make an output current 266 constant and provide the output current 266 to the one or more LEDs 250, by working with at least the capacitor 238.\nAs shown in FIG. 2, a drain voltage at a drain terminal of the transistor 262 is received by a voltage divider that includes the resistors 276 and 278. The drain terminal of the transistor 262 is connected to the winding terminal 265 of the winding 260, and the winding terminal 263 of the winding 260 is connected to the capacitor 230 and the resistor 279. In response, the voltage divider generates a voltage 277, which is received by the controller terminal 256. The PWM controller 240 uses the voltage 277 to detect the end of a demagnetization process of the winding 260. The detection of the end of the demagnetization process is used to control an internal error amplifier of the PWM controller 240, and through the controller terminal 246, to control charging and discharging of the capacitor 234.\nAlso, after the PWM controller 240 has started up, the resistor 274 is used to detect a current 261, which flows through the winding 260. The current 261 flows from the winding 260 through the resistor 274, which in response generates a sensing voltage 275. The sensing voltage 275 is received by the PWM controller 240 at the controller terminal 258, and is processed by the PWM controller 240 on a cycle-by-cycle basis. The peak magnitude of the sensing voltage 275 is sampled, and the sampled signal is sent to an input terminal of the internal error amplifier of the PWM controller 240. The other input terminal of the internal error amplifier receives a reference voltage Vref.\nAs shown in FIG. 2, the rectified output voltage 222 is received by a voltage divider that includes the resistors 270 and 272. In response, the voltage divider generates a voltage 271, which is received by the controller terminal 242. The PWM controller 240 processes the voltage 271 and determines phase angle of the voltage 271. Based on the detected range of phase angle of the voltage 271, the PWM controller 240 adjusts the reference voltage Vref, which is received by the internal error amplifier.\nThe bleeder 224 is used to ensure that, when the TRIAC dimmer 210 is fired on, an input current 264 that flows through the TRIAC dimmer 210 is larger than a holding current required by the TRIAC dimmer 210, in order to avoid misfire of the TRIAC dimmer 210 and also avoid flickering of the one or more LEDs 250. For example, the bleeder 224 includes a resistor, which receives the rectified output voltage 222 at one resistor terminal of the resistor and receives the ground voltage at the other resistor terminal of the resistor. The resistor of the bleeder 224 allows a bleeder current 268 to flow through as at least part of the input current 264. In another example, if the holding current required by the TRIAC dimmer 210 is small and if the average current that flows through the transistor 262 can satisfy the holding current requirement of the TRIAC dimmer 210, the bleeder 224 is not activated or is simply removed.\nAs shown in FIG. 2, the lighting system 200 includes, for example, a quasi-resonant system with a buck-boost topology. The output current 266 of the quasi-resonant system is received by the one or more LEDs 250 and is determined as follows:\n                              I          o                =                              1            2                    \u00d7                                    V              ref                                      R              cs                                                          (                  Equation          \u2062                                          \u2062          1                )            where I0 represents the output current 266 of the quasi-resonant system of the lighting system 200. Additionally, Vref represents the reference voltage received by the internal error amplifier of the PWM controller 240. Moreover, Rcs represents the resistance of the resistor 274.\nFIG. 3 is a simplified diagram showing certain conventional components of the lighting system 200 as shown in FIG. 2. The pulse-width-modulation (PWM) controller 240 includes a dimming control component 300 and a transistor 350. The dimming control component 300 includes a phase detector 310, a reference voltage generator 320, a pulse-width-modulation (PWM) signal generator 330, and a driver 340.\nFIG. 4 shows certain conventional timing diagrams for the lighting system 200 as shown in FIGS. 2 and 3. The waveform 471 represents the voltage 271 as a function of time, the waveform 412 represents the phase signal 312 as a function of time, the waveform 475 represents the sensing voltage 275 as a function of time, and the waveform 464 represents cycle-by-cycle average of the input current 264 as a function of time.\nAs shown by FIGS. 3 and 4, the lighting system 200 uses a closed loop to perform dimming control. The phase detector 310 receives the voltage 271 through the terminal 242, detects phase angle of the voltage 271, and generates a phase signal 312 that indicates the detected range of phase angle of the voltage 271. As shown by the waveform 471, the voltage 271 becomes larger than a dim-on threshold voltage (e.g., Vth_dimon) at time ta and becomes smaller than a dim-off threshold voltage (e.g., Vth_dimon) at time tb. The dim-on threshold voltage (e.g., Vth_dimon) is equal to or different from the dim-off threshold voltage (e.g., Vth_dimoff). The time duration from time ta to time tb is represented by TR, during which the phase signal 312 is at the logic high level, as shown by the waveform 412. The time duration TR represents the detected range of phase angle of the voltage 271.\nDuring the time duration TR, the sensing voltage 275 ramps up and down. For example, during the time duration TR, within a switching period (e.g., TSW), the sensing voltage 275 ramps up, ramps down, and then remains constant (e.g., remains equal to zero) until the end of the switching period (e.g., until the end of TSW).\nThe phase signal 312 is received by the reference voltage generator 320, which uses the detected range of phase angle of the voltage 271 to generate the reference voltage 322 (e.g., Vref). As shown in FIG. 3, the reference voltage 322 (e.g., Vref) is received by the PWM signal generator 330. For example, the PWM signal generator 330 includes the internal error amplifier of the PWM controller 240. In another example, the PWM signal generator 330 also receives the sensing voltage 275 and generates a pulse-width-modulation (PWM) signal 332. The PWM signal 332 is received by the driver 340, which in response generates a drive signal 342 and outputs the drive signal 342 to the transistor 350. The transistor 350 includes a gate terminal, a drain terminal, and a source terminal. The gate terminal of the transistor 350 receives the drive signal 342. The drain terminal of the transistor 350 is coupled to the controller terminal 254, and the source terminal of the transistor 350 is coupled to the controller terminal 258.\nAs shown by the waveform 475, the reference voltage 322 (e.g., Vref) is used by the PWM signal generator 330 to generate the PWM signal 332, which is then used to control the peak magnitude (e.g., CS_peak) of the sensing voltage 275 for each PWM cycle during the time duration TR. For example, each PWM cycle corresponds to a time duration that is equal to the switching period (e.g., TSW) in magnitude. In another example, if the detected range of phase angle of the voltage 271 (e.g., corresponding to TR) becomes larger, the reference voltage 322 (e.g., Vref) also becomes larger. In yet another example, if the detected range of phase angle of the voltage 271 (e.g., corresponding to TR) becomes smaller, the reference voltage 322 (e.g., Vref) also becomes smaller.\nAccording to Equation 1, if the reference voltage 322 (e.g., Vref) becomes larger, the output current 266 (e.g., Io) of the quasi-resonant system of the lighting system 200 also becomes larger; if the reference voltage 322 (e.g., Vref) becomes smaller, the output current 266 (e.g., Io) of the quasi-resonant system of the lighting system 200 also becomes smaller.\nAs shown by FIG. 2, the cycle-by-cycle average of the input current 264 is approximately equal to the sum of cycle-by-cycle average of the output current 266 (e.g., Io) and the bleeder current 268. During the time duration TR, within each switching cycle of the PWM signal 332, the output current 266 changes with time, so the average of the output current 266 within each switching cycle is used to determine the cycle-by-cycle average (e.g., I_PWM_av) of the output current 266 as a function of time. When the time duration TR becomes smaller, the reference voltage 322 (e.g., Vref) also becomes smaller and the one or more LEDs 250 are expected to become dimmer. When the time duration TR becomes too small, the reference voltage 322 (e.g., Vref) also becomes too small and the cycle-by-cycle average (e.g., I_PWM_av) of the output current 266 during the time duration TR becomes smaller than the holding current (e.g., I_holding) required by the TRIAC dimmer 210. In order to avoid misfire of the TRIAC dimmer 210 and also avoid flickering of the one or more LEDs 250, the bleeder current 268 (e.g., I_bleed) is provided in order to increase the cycle-by-cycle average of the input current 264 during the time duration TR. As shown by the waveform 464, the cycle-by-cycle average of the input current 264 during the time duration TR becomes larger than the holding current required by the TRIAC dimmer 210.\nAs shown in FIG. 3, the driver 340 outputs the drive signal 342 to the transistor 350. The transistor 350 is turned on if the drive signal 342 is at a logic high level, and the transistor 350 is turned off if the drive signal 342 is at a logic low level. When the transistor 262 and the transistor 350 are turned on, the current 261 flows through the winding 260, the transistor 262, the controller terminal 254, the transistor 350, the controller terminal 258, and the resistor 274. If the transistor 350 becomes turned off when the transistor 262 is still turned on, the transistor 262 then also becomes turned off and the winding 260 starts to discharge. If the transistor 350 becomes turned on when the transistor 262 is still turned off, the transistor 262 then also becomes turned on and the winding 260 starts to charge.\nAs shown in FIGS. 2-4, the lighting system 200 uses a closed loop to perform dimming control. For example, the lighting system 200 detects the range of phase angle of the voltage 271, and based on the detected range of phase angle, adjusts the reference voltage Vref that is received by the internal error amplifier of the PWM controller 240. In another example, the lighting system 200 provides energy to the one or more LEDs 250 throughout the entire time period of each switching cycle during the time duration TR, which corresponds to the unclipped part of the waveform of the AC input voltage 214 (e.g., VAC).\nAs discussed above, a bleeder (e.g., the bleeder 224) can help a lighting system (e.g., the lighting system 200) to meet the holding-current requirement of a TRIAC dimmer (e.g., the TRIAC dimmer 210) in order to avoid misfire of the TRIAC dimmer (e.g., the TRIAC dimmer 210) and avoid flickering of one or more LEDs (e.g., the one or more LEDs 250). But the bleeder (e.g., the bleeder 224) usually increases heat generation and reduces energy efficiency of the lighting system (e.g., the lighting system 200). Such reduction in energy efficiency usually becomes more severe if a bleeder current (e.g., the bleeder current 268) becomes larger. This reduced energy efficiency often prevents the lighting system (e.g., the lighting system 200) from taking full advantage of high energy efficiency and long lifetime of the one or more LEDs (e.g., the one or more LEDs 250).\nHence it is highly desirable to improve the techniques of dimming control."}
-{"text": "Field of the Invention\nEmbodiments of the present invention generally relate to performing capacitance sensing while updating a display, or more specifically, to performing capacitance sensing when display updating is paused.\nDescription of Related Art\nInput devices including proximity sensor devices (also commonly called touchpads or touch sensor devices) are widely used in a variety of electronic systems. A proximity sensor device typically includes a sensing region, often demarked by a surface, in which the proximity sensor device determines the presence, location and/or motion of one or more input objects. Proximity sensor devices may be used to provide interfaces for the electronic system. For example, proximity sensor devices are often used as input devices for larger computing systems (such as opaque touchpads integrated in, or peripheral to, notebook or desktop computers). Proximity sensor devices are also often used in smaller computing systems (such as touch screens integrated in cellular phones)."}
-{"text": "DE 10 2012 010 757 A1 discloses an illuminating device for a vehicle. The illuminating device includes an eye-tracking system for detecting an eye position and its viewing direction. The illuminating device is divided into several predefined switch-on ranges, and light is emitted only into that predefined switch-on range which corresponds to the user's field of vision. No light is emitted as soon as the user's field of vision does not correspond to the respective predefined switch-on range.\nThe disadvantage, however, consists in that drivers of a motor vehicle during darkness frequently orientate themselves using the lighting of motor vehicle instruments. If the instruments are always completely switched-off as soon as the driver ceases to look in their direction, the driver loses his orientation more easily, which may make operation of the motor vehicle more difficult.\nThis leads to the requirement to further develop a method to address this disadvantage in such a way that it becomes easier for the driver to orientate himself in a dark motor vehicle without being irritated by the lighting of at least one motor vehicle instrument."}
-{"text": "Objective lens assemblies are commonly used in microscopes, telescopes, cameras and other devices for gathering light from an object being observed and focusing the light to form an image of the object. Objective lens assemblies that operate in visible spectrum of light are quite common.\nCurrently, the applicant of the present invention is developing a microscope that operates in the mid infrared (\u201cMIR\u201d) light spectrum. Unfortunately, existing objective lens assemblies do not provide sufficient performance in the MIR light spectrum."}
-{"text": "The present invention relates to connection members for components of a close-coupled pressurized system and, more particularly, a connector spool assembly provided with adjustment components to allow movement of the connector spool to facilitate separation and removal of system components.\nOne type of compression system is a compressor close-coupled to an electric motor driver, which provides for a compact design with significant benefits over traditional base-plate mounted compressor trains. A motor casing and a compressor casing comprise separate bodies requiring removal for service. One problem with component removal service activity is the cost and time required to disconnect process piping and instrumentation connected to each casing. Individual case removal is especially problematic for applications where the unit has compressor casings at each end of a double ended motor drive."}
-{"text": "The present disclosure relates to a display unit displaying an image, and an image processing unit for use in such a display unit, and a display method.\nRecently, a cathode ray tube (CRT) display unit has been actively replaced with a liquid crystal display unit or an organic electro-luminescence (EL) display unit. The liquid crystal display unit and the organic electro-luminescence display unit are each being a mainstream display unit due to low power consumption and a flat configuration thereof.\nDisplay units are in general desired to have high image quality. Image quality is determined by various factors including contrast. Increase of peak luminance may be a technique for improving contrast. Specifically, reduction of a black level is limited by reflection of outside light, etc. Hence, in the above technique, peak luminance is increased (extended) to improve contrast. For example, Japanese Unexamined. Patent Application Publication No. 2008-158401 (JP-A-2008-158401) discloses a display unit, in which an increasing level (extending level) of peak luminance and gamma characteristics are each varied depending on an average of image signals to achieve improvement in image quality and reduction in power consumption.\nIn some display units, each pixel is configured of four sub-pixels. For example, Japanese Unexamined Patent Application Publication No. 2010-33009 discloses a display unit, in which each pixel is configured of sub-pixels of red, green, blue, and white to improve luminance or reduce power consumption, for example."}
-{"text": "The invention relates to a dispenser or an assembly suitable as a dispenser, serving as a receptacle, reservoir and/or discharger for media which may be liquid, pasty, powdery and/or gaseous. All components of the dispenser or assembly may be made of plastics or as compression or injection molded components. For discharge the dispenser can be freely held and simultaneously actuated single-handedly. Its length thus amounts to max. 10 cm or 7 cm, its largest width not more than 8 cm or 5 cm. The dispenser is suitable for dispensing single droplets of the medium, a jet or atomized particle or droplet aerosol thereof. Furthermore, the dispenser may be configured for discharging but a single dose of the medium or for a single stroke with no return stroke or for repeated discharges each with a spring-actuated return stroke inbetween.\nExperience has shown it to be expedient to compose complicated assemblies of components molded separately which during molding are located with or without a direct joint spaced away from each other or in a position other than that required in the operating condition. Reference is made to the German laid-open document 196 05 153 as well as to the pending German patent 198 13 078.3 in including the features and effects described therein in the present invention.\nThe invention is based on the object of providing a dispenser or a method of producing an assembly for a dispenser or the like which avoids the disadvantages of known configurations. It is more particularly the object to provide assemblies which have an increasing or decreasing inner or outer cross-section in the opposing direction. The dispenser is intended for facilitated production and safe operation.\nIn accordance with the invention two or more components are produced at the same time or with the same flow of plasticated material, immediately demolded once solidified or released in some other way at their jointing zones and then directly positioned relative to each other so that they can then be combined into an assembly. For the components the same material or differing materials may be employed. The components pass through the same temperature curves at the same time up to solidification and may have the same or differing volume of material. Expediently the components are produced in the same mold or so that they adjoin one or more common parts of the mold each integrally. This applies more particularly to the jointing surface areas of the components moldable juxtaposed in common by a movable part of the mold. After the components have solidified and subsequent retraction thereof or of another part of the mold these jointing zones are located exposed. The components can then be moved relative to each other until joined together and demolded completely where necessary. It is good practice when the components are located in production axially parallel or directly juxtaposed almost in contact with each other. Once the one component has been joined to the other it forms an elongation of the other component in the direction of its greatest extent. After being joined, forming the operating condition of the assembly for operation of the dispenser, the two components merge into a length which is smaller than the length of the one or other component. The components may, however, also be face joined without any mutual longitudinal engagement and locked in place mutually by a further component. Thus mutual locking of the components may be with zero clearance or positive, namely by being radially centered or by a captive lock.\nAlthough the configuration in accordance with the invention is suitable for the outer or base bodies of dispensers it is particularly expedient for core bodies. One such core body is located totally concealed in the interior of the dispenser of the corresponding base body, e.g. within a discharge nozzle. This base body may also form the third component for the cited locational lock. Advantageously one or both components of the assembly forms longitudinally a middle section of largest outer width, w adjoining at each end thereof an end section of comparitively reduced outer width. Each of the end sections is formed by another component. An end section may be a hollow needle having a smallest diameter at the tip of the needle of less than one millimeter and a length of less than 10 or 8 mm. The other end section may be a dished, fluted or outer face-recessed body having a radially protruding collar forming the shorter longitudinal part of the middle section.\nThe two components are advantageously joined to each other via a single connecting member or link directly joining each of the components by a link section. The link sections are then mutually movable and adjoined by a connecting location which may remain stationary in mutual movement of the components such as in movement of the corresponding link section relative to the corresponding component. The connecting location is expediently a hinging zone having a sole hinging axis and/or a designed frangible location at which the link sections are parted in mutual movement of the components and prior to attaining the operating position in forming opposing fractured surface areas. The mold cavity for the link may form the one or sole flow channel via which the plasticated material flows from the mold cavity for the one component, more particularly the larger volume component, into the mold cavity for the other component. The smallest cross-section of this channel and thus of the link may be less than 5, 2 or one tenth of a mm2.\nIn production the jointing surface areas of the two components later to directly adjoin in the operating position are expediently located in the same plane. Beyond one of these jointing surface areas a locking member or the like may protrude. In production these jointing surface areas may point in the same direction or in opposite directions. Up to each jointing surface area the link may also extend which may comprise a surface directly translating into the jointing surface areas in the same plane or frangible or parting surface areas in this plane after parting. The components may also be translated by a radial or linear movement into their operating position, the one component forming a sliding guide for the other component flanked only at the bottom and sides which, however, does not attain the guide until after a first portion of the shifting travel or after the link has been parted. In addition, the components may be produced separately and then assembled in accordance with the invention.\nIrrespective of the configuration as described, the dispenser is configured more particularly as a receptacle and reservoir for biological active substances over several weeks, months or even years. These may be physiological active substances containing hormones and/or cleavage products such as peptides containing protein. Such biological information transmitters which may contain amino acids and other similar active substances may be highly sensitive to moisture, this being the reason why they are held in the dispenser in a pressure-tight chamber which is not opened until immediately prior to delivery from the dispenser, e.g. by a closure being ruptured by means of the cited assembly.\nThese and further features of the invention also read from the description and the drawings, each of the individual features being achieved by themselves or severally in the form of sub-combinations in one embodiment of the invention and in other fields and may represent advantageous aspects as well as being patentable in their own right, for which protection is sought in the present."}
-{"text": "In current technologies, the threshold voltage of semiconductor devices does not scale with the power supply voltage and ground rules because of the non-scalability of the sub-threshold slope. Thus, the minimum gate oxide thickness and/or maximum wordline boost voltage of the array MOSFET is constrained by reliability considerations.\nWhen used for the support MOSFET, the relatively thick gate oxide (having a thickness of greater than \u22486 nm for deep sub-\u03bcm technology) required by the array MOSFET results in degradation in the performance of the MOSFET device. Furthermore, if a thinner gate oxide is used to improve the performance of the support circuitry, charge transfer efficiency in the device array is compromised as a result of the reliability limitation of the wordline boost voltage.\nIdeally, in such technology, a dual gate oxide thickness is desired. In the prior art, it is known to subject the array transistor to a dual gate oxidation process or an alternative gate oxidation process as compared to the support circuitry. These additional gate oxidation processing steps are costly, and they are also yield limiting since one must utilize additional processing steps such as, but not limited to: masking, exposure, etching, oxidizing and strip masking, which grow a second oxide on the entire structure of the MOSFET device. As such, prior art gate oxidation processes are not reliable nor cost efficient.\nIn view of the drawbacks mentioned above with prior art processes of fabricating MOSFETs, there is a continued need for providing a new and improved method of fabricating a MOSFET and other devices in which a dielectric layer, e.g., gate oxide, having a dual thickness can be formed without adding extra processing steps and costs to the overall manufacturing process."}
-{"text": "A variety of different methods have been developed to assay oligonucleotides, including DNA or RNA fragments. Such assays are typically directed to determining whether a sample includes oligonucleotides having a particular target oligonucleotide sequence. In some instances, oligonucleotide sequences differ by only a few nucleotides, as in the case of many allelic sequences. Single nucleotide polymorphisms (SNPs) refer to alleles that differ by a single nucleotide. Even this single nucleotide difference can, at least in some instances, change the associated genetic response or traits. Accordingly, to determine which allele is present in a sample, the assay technique must be sufficiently sensitive to distinguish between closely related sequences.\nMany assay techniques include multiple components, each of which hybridizes to other component(s) in the assay. Non-specific hybridization between components (i.e., the hybridization of two non-complementary sequences) produces background noise in the assay. For example, closely related, but not identical, sequences can form imperfect duplexes in which base pairing is interrupted at positions where the two single strands are not complementary. Non-specific hybridization increases when the hybridizing components have similar sequences, as would be the case, for example, for many alleles and particularly for SNP alleles. Thus, for example, hybridization assays to determine which allele is present in a sample would benefit from methods that reduce non-specific hybridization or reduce the impact of non-specific hybridization on the assay."}
-{"text": "A. Field of Invention\nThis invention pertains to a method and apparatus for operating an internal combustion engine using a fuel consisting of water and a water-soluble flammable substance that is injected into a mixture of hydrogen and air.\nB. Description of the Prior Art\nThe use of fossil fuels to run engines that used, for example, in cars and other vehicles, as well as many other engines used for a variety of purposes, is based on a very old concept based on the internal combustion engines developed in the nineteenth century. Despite intense research and development for alternate fuels for the last 50 years, fossil fuel derived from petroleum or natural gas, is still essentially the primary source of energy almost all the internal combustion engines presently in use all over the world.\nAs a result, the world supply of fossil fuels have been severely depleted creating a shortage, and the price of oil has been climbing for the past 40 years. In addition such fuels are very polluting and some suggest that it has either been the primary cause or has contributed substantially to global warming. All these factors led to many efforts to find and harness renewable energy sources other than traditional fossil fuels. Several alternative fuels have been introduced in the past few years to reduce the impact of petroleum depletion, including hybrid cars, electric cars, bio diesel, hydrogen based cars, etc. However, none of these solutions were effective. One reason for this lack of success is that they require a completely new infrastructure for the production of the engines, as well as the production and distribution of the fuel. Moreover, the most solutions proposed so far were incompatible with the existing engines and, therefore. The cost of replacing all the existing fossil burning engines may be so high that it may render any solution based on alternate fuels unacceptable, at least, in a short term basis.\nWater as a source of fuel has been suggested by many in the past and many experiments have been conducted testing such systems. The basis of such experiments is the fact that water can be separated in to hydrogen and oxygen and the resulting stoichiometric mixture can be fed in to an internal combustion engine to generate power. However past experiments yielded unsatisfactory results. The main obstacle for their success is based on the fact that the energy required to separate the water into its components is much greater than the energy produce by the engine. In addition the amount H2 mixture needed to run a typical automotive engine is too large to make such a system practical.\nSystems are presently available on market that can be used as accessories or add-ons to internal combustion engines using fossil fuels, however independent tests have shown that, in fact, these systems have very little, if any, effect on the overall efficiency of the engine.\nA system developed by the present inventors is described in two co-pending applications includes means of generating from water and supplying a small amount of hydrogen/oxygen gas mixture into a standard internal combustion engine. (See U.S. Patent Application Publications 2010/0122902 and 20110203917). More specifically, these co-pending applications describe an efficient process and apparatus for generating a two-to-one mixture of hydrogen and oxygen, commonly referred to a brown gas or HHO. The mixture helps increase the efficiency of the conventional internal combustion engine by burning the fossil fuel more efficiently. While this latter system is much more efficient that previously described systems; its efficiency is still limited by the amount of hydrogen and oxygen produced on board a vehicle. Moreover, the internal combustion engine described is still burning a fossil fuel."}
-{"text": "With the popularity of digital photography and digital image processing, consumers have increasingly desired to transfer photographic images stored on conventional film negatives into electronically stored digital images. Typically, this is accomplished by loading a sheet of processed film into a scanner and scanning the film to produce the digital image. Processed film is normally cut into sheets containing one to six images. Thus, if a user has a large number of negatives to scan, the process of loading each individual sheet of film into the scanner can become overly time-consuming. Accordingly, there is a desire for improved systems and methods for automate the loading and scanning of multiple sheets of film.\nConventional systems for handling the feeding of paper or film documents, such as those used in photocopiers, printing presses, printers and scanners, are not well suited for the handling of film. In particular, the rollers used for feeding individual sheets from a stack of paper or film may damage the image on sheets of film. In addition, these loading mechanisms are configured to load a large number of identically-sized sheets of paper in standard sizes such as 8.5\u2033\u00d711\u2033 or 8.5\u2033\u00d714\u2033. In contrast, photographic film negatives are often manually cut, resulting in sheets of film of varying lengths that are difficult to accurately load on a bulk basis. In addition, photographic film can change its shape over time or during operation, such as when the film curls around unpredictable angles."}
-{"text": "Some homes today are equipped with smart home networks to provide automated control of devices, appliances and systems, such as heating, ventilation, and air conditioning (\u201cHVAC\u201d) system, lighting systems, alarm systems, home theater and entertainment systems. Smart home networks may include control panels that a person may use to input settings, preferences, and scheduling information that the smart home network uses to provide automated control the various devices, appliances and systems in the home. For example, a person may input a desired temperature and a schedule indicating when the person is away from home. The home automation system uses this information to control the HVAC system to heat or cool the home to the desired temperature when the person is home, and to conserve energy by turning off power-consuming components of the HVAC system when the person is away from the home. Also, for example, a person may input a preferred nighttime lighting scheme for watching television. In response, when the person turns on the television at nighttime, the home automation system automatically adjusts the lighting in the room to the preferred scheme."}
-{"text": "1. Field of the Invention\nThe present invention relates to a garnet-type ion conducting oxide, a complex, a lithium secondary battery, a manufacturing method of a garnet-type ion conducting oxide and a manufacturing method of a complex.\n2. Description of the Related Art\nGarget-type oxides such as Li7La3Zr2O12 and Li7ALa3Nb2O12 (A=Ca, Sr or Ba) synthesized by the solid-phase reaction method have been proposed conventionally as a solid electrolyte configured to conduct lithium ion (Non-Patent Literatures 1 to 3). It has been reported that this solid electrolyte has the conductivity of 1.9 to 2.3\u00d710\u22124 Scm\u22121 (25\u00b0 C.) and activation energy of 0.34 eV. The inventors have studied a solid electrolyte of Li7La3Zr2O12-based garnet-type ion conducting oxide among garnet-type oxides having excellent chemical stability and a wide potential window. For example, it has been proposed that the Zr sites in this solid electrolyte should be substituted with an element such as Nb, in order to enhance the conductivity (see, for example, Patent Literature 1). This solid electrolyte has high conductivity but needs treatment at high temperature such as 1200\u00b0 C. It has been proposed, on the other hand, that La sites should be additionally substituted with an alkaline-earth metal, in order to minimize reduction of the electric conductivity and reduce the firing energy (see, for example, Patent Literature 2).\nA solid electrolyte including Li, La, Zr, O and Al has been proposed as Li7La3Zr2O12-based solid electrolyte (see, for example, Patent Literature 3). According to the disclosure of this prior art, addition of Al to Li7La3Zr2O12-based solid electrolyte provides the solid electrolyte with the density and the conductivity required for the solid electrolyte material."}
-{"text": "1. Field of the Invention\nThe present invention relates to the test techniques and testability architectures used in testing integrated circuits. More specifically, the invention pertains to test techniques applied to testing of user-configurable arrays before they are configured by the user.\n2. The Prior Art\nUser-configurable gate arrays consist of logic circuits or blocks that can be connected together by configurable interconnections, such as anti-fuse elements, to implement a desired circuit function. The configurable interconnect consists of interconnect layers such as metallization, and configurable devices which, when programmed, establish electrical connections between the interconnect layers. However, before configuring the circuit to implement a particular function, all the individual modules in the array and all the input/output (I/O) modules and buffers are isolated from one another. This presents a challenging test problem.\nBefore the circuit is configured by the user, all of the active circuits in such integrated circuits such as logic modules, I/O modules, configuring circuits, etc., must be tested and guaranteed to be fully functional and meet all required specifications. In addition, all passive interconnect circuits such as metallization interconnect, anti-fuse elements, feed-thru pass transistors, must also be free of defects and guaranteed. This is necessary so that a customer configuring such a circuit can expect a fully functional, high quality integrated circuit after his application circuit is mapped into the device. It is thus imperative that test architectures and test techniques be developed to solve this problem, namely, how to guarantee full functionality and spec of a one-time programmable user configurable array circuit before the circuit is configured by the user.\nUser configurable arrays or PLDs (programmable logic devices) which use erasable elements to implement their interconnect do not have to contend with this problem since the array can be configured to implement any circuit pattern, be fully tested and later erased to the \"blank\" state for reconfiguration."}
-{"text": "1. Field of the Invention\nThe invention relates generally to presses and, more particularly, to shell presses and associated methods for forming container closures or ends, commonly referred to as shells. The invention also relates to die assemblies for shell presses.\n2. Background Information\nThe forming of can ends or shells for can bodies, namely aluminum or steel cans, is generally well-known in the art.\nThere is an ongoing desire in the can-making industry to manufacture shells as rapidly and efficiently as possible. Among the ways companies have attempted to achieve these objectives are: (1) to increase the number of pockets in the die set, within which shells can be formed; and (2) to increase the speed (e.g., strokes per minute (spm)) at which the shell press operates. In general, with each stroke of the shell press ram, one shell is formed in each tooling pocket of the die assembly. Thus, a 24-out die assembly, for example, which has 24 tooling pockets, is capable of forming 24 shells, per stroke. U.S. Pat. No. 5,491,995, which is hereby incorporated herein by reference, discloses an example of a relatively high capacity (e.g., without limitation, operating speed of up to 400 spm, or more) end shell manufacturing system having a 24-out die assembly.\nHowever, forming shells at relative high speeds generates heat. The heat, which is caused by the friction associated with drawing the metal over forming surfaces of the die assembly and/or clamping the metal between various pressure pads and drawing it through reduced tooling clearances to provide a desired shape, can be excessive, resulting in thermal expansion of the die shoes. Among other disadvantages, such thermal expansion undesirably shifts tooling and/or reduces critical clearances between cutting and/or forming tools. Consequently, tooling wear or damage can result and/or certain features of the end shells are manufactured out-of-specification. For example, thinned spots can be created in the material from which the end shell is manufactured, leading to a loss in buckle pressure performance in the final product.\nThe foregoing difficulties have been exacerbated by the development of new shell designs having aggressive material thicknesses and shapes. For example, some shells require reduced material thickness and/or have a relatively complex geometry. Such shapes often necessitate additional pressure pads and increased forming pressures in order to properly manufacture the end shells.\nPrior proposals that attempted to address thermal expansion of the die assembly tooling (e.g., without limitation, upper and lower die shoes) involved aligning the upper tooling with respect to the lower tooling in the die assembly in a manner intended to compensate for the thermal expansion. Other proposals require coolant (e.g., chilled water) to be pumped throughout the die assembly, for example, to reduce the rate and amount of thermal expansion of the die shoes. However, estimating and establishing the proper aligning of the upper tooling with respect to the lower tooling is a time-consuming process, and it can be difficult to maintain the desired alignment. Similarly, systems that add coolant or other suitable additional cooling or heating mechanisms to the die assembly to compensate for thermal expansion, are costly to install and maintain.\nThere is, therefore, room for improvement in shell presses, and in die assemblies and associated methods therefor."}
-{"text": "1. Field of the Invention\nThe present invention pertains to an eyeglass end face machining method, particularly to the polishing to a mirror polishing that is performed on the end face after bevel edging, or the planing, such as machining to a mirror polishing, that is performed on the end face after edging.\n2. Description of the Related Art\nThe lens end face of rimless eyeglasses lenses usually referred to as three-piece eyeglass lenses is exposed and not covered by a rim, etc., and therefore, they must have a surface that has been polished until glossy. In response to this need, technology has been presented whereby eyeglass lenses, whose end face has thus far been smoothed manually in order to obtain a face that has been polished until glossy, are mechanically polished by placing a movement mechanism with tracing capability in the polishing wheel part (for instance, Japanese Patent Laid-Open No. Sho 64-87144). This grinds inclined faces, such as the end face of polyhedron cut lenses, etc., and although the shape around the eyeglass lens is complex because of the polyhedron cut, the end face itself, which becomes the surface to be grounded, is a flat surface and simple. Consequently, the above-mentioned technology cannot be used when the surface to be polished itself has a complex shape, such as lens end faces with a bevel. Now, because the lens end face with a bevel is usually concealed by the rim of the frame and there is no need to polish the bevel faces, a lens end face with a bevel itself is usually not polished.\nHowever, there has been a demand in recent years for thin rims in order to obtain frames that are more lightweight and fashionable, etc., and it is often the case that if the lens fitted into the rim is a strong-minus-power lens with a thick edge, the lens will protrude from the rim of the frame. It is pointed out that the bevel faces remains white when polishing of the lens end face is completed by bevel-polishing and this poses a problem aesthetically. Polishing the bevel surface that remains white until it is transparent is only accomplished by buff polishing the bevel surface by hand, etc., and this takes time and increases cost.\nThe objective of the present invention is to solve the above-mentioned problems with prior art by mechanically polishing the bevel faces in 2 steps and to present a lens end face machining method, wheel and device for eyeglass lens end face machining with which it is possible to speed up the polishing process and make finishing precision uniform and obtain fashionable eyeglass.\nMoreover, in addition to the aesthetic problem of the lens end face remaining white after bevel polishing that was previously described, there is a problem with polishing precision and fashionable eyeglasses in that when planing, such as smooth machining and machining to a mirror polishing etc., is performed with a wheel that has a bevel-groove and a planing face, streaks are made. That is, cylindrical grinding stone called diamond wheels(stone) have a bevel-groove for formation of a bevel at the end face of the eyeglass lens and a flat face for flat machining the end face of an eyeglass lens. In further detail, the wheel has groove inclined face 301 for V finishing having a specific angle with respect to the axial direction called angle No. 1, flank 203 for the eyebrow of the frames continuous with this groove inclined face 301 having a specific angle with respect to the axial direction referred to as angle No. 2 that is smaller than angle No. 1, and flat finishing face 303 continuous with this flank 302 for flat machining and parallel to the axial direction on the surface around the periphery of the wheel. The inclination at boundary K between above-mentioned flank 302 and flat finishing face 303 is not continuous.\nConsequently, when an eyeglass lens moves past boundary K to the left in the direction of the X axis during flat machining, apex A of the end face of eyeglass lens 6 straddles boundary K and a streak from boundary K is made in end face 6a of eyeglass lens 6. When a streak is made in end face 6a of the eyeglass lens, edging precision drops and becomes non-uniform, and the product is not fashionable. Therefore, such a streak is undesirable. This is particularly a problem with flat finished surfaces that remain white and are further given a mirror finish so that they are transparent.\nThereupon, in order to solve this problem, flat finishing face 303 is made longer in the axial direction so that even if eyeglass lens 6 moves to the left in the direction of the X axis during flat edging, it will not pass boundary K. However, there is a problem in that as a result, wheel 1 is larger.\nInerdentally there is a demand for mechanical polishing of the bevel face that remains white using a wheel as a means of solving the above-mentioned aesthetic problem of the lens end face remaining white after bevel polishing because buffing, etc., manually takes time and increases cost. However, there is also a problem when a polishing wheel is used with the existing wheel in that the device becomes bigger.\nThe objective of the present invention is to solve the above-mentioned problems of prior art and present an eyeglass lens end face matching method with which polishing precision is uniform, the product is excellent in terms of being fashionable, and the device can be reduced in size. Another objective of the present invention is to present an eyeglass end face machining method with which it is possible to add a polishing wheel that can give the eyeglass lens end face a mirror polish without greatly increasing length of the wheel in the axial direction."}
-{"text": "With the proliferation of the automated interactive machines, exemplified by the automated teller machines (ATM) for financial transactions, there has been an emerging need for a more reliable personal identification system for authenticating users who desire to conduct transactions remotely and automatically without human intervention. Conventionally, a person simply inserts her ATM card into the machine to have her account information and password, or PIN (\"Personal Identification Number\", used here interchangeably with the word \"Password\"), read. However, as the everyday life as a whole becomes more automated and security-conscious, a person often has to manage various different passwords and PIN's, for accesses to her banking account, her home security system, or her eMail account, to name just a few. This overflow of information has already contributed to the complexity of conventional personal identification systems in that without the correct password for an ATM, a legitimate user may be denied of her access to her account or her on-line brokerage account.\nThere is an often overlooked burden placed on the institutions providing on-line, or remote, transactions which are accessed through the customers' passwords or PINs. Maintaining the passwords or PINs forces the financial institutions to allocate additional machines and human resources to manage interface with customers when a customer forgets her Pin or when a customer requests her PIN be changed.\nAlso, passwords have been proven to be insufficient in preventing fraud, where all a would-be criminal needs is an ATM card and the password, which are both reasonably within the reach of those unscrupulous ones. This is just the first example of how the conventional personal identification paradigm is vulnerable, in addition to being complex as discussed above.\nAnother problem plagues the integrity of the supposedly secured financial transaction, where sometimes it is the actual account holder who defrauds the institution by first accessing her account and later denying such transaction from ever taking place. While there is a limit as to the extent of this sort of heinous behavior, it amounts to a significant sum even with just a small percentage of the ATM transactions considered. Without a more reliable identification system, institutions will just have to write off the losses or pass the losses to the rest of the consumers, thereby increasing everyone's cost of doing business.\nAside from the ATM transactions, with the increasing affordability, as well as sophistication, of personal computers and telecommunication hardware and software, it is more likely that one will soon be accessing a host of information or conducting a variety Of secured transactions using a PC, a modem and a common public switching network, such as Prodigy and Internet, etc. Authentication thus becomes an even more paramount task for the industry to tackle.\nA simple personal identification system may address the above problems. Fingerprints have been known years ago to have a high degree of accuracy and reliability. One never forgets her fingerprints, or confuses the fingerprints with other information. Also, a criminal cannot steal or duplicate someone's fingerprints to impersonate the account holder, generally speaking. Therefore, fingerprints are essentially a personal identification with a one-to-one correspondence, given that the fingerprint recognition systems have progressed along with the information revolution. Companies such as Identix and Startech have developed front-end fingerprint image recognition systems to reliably and accurately analyze and recognize fingerprints.\nAt the back-end, major processor suppliers such as IBM and AT&T already have systems in place to provide a linkup with the fingerprint image recognition systems such that the massive fingerprint database may be linked and accessed for the institution to quickly authenticate the person in front of its machine, or the person seeking to access her brokerage account through a PC with a modem. To a certain extent, the present front-end and back-end suppliers have reached a point, where it is merely a matter of time before their capabilities and achievements can be fully utilized by the industry, especially the financial industry.\nEven with reliable fingerprint image recognition systems at the front-end and quick-response processor at the back-end, there are still problems with this paradigm. Assuming it is reasonably affordable for a PC owner to have a personal fingerprint recognition device to provide access to her on-line brokerage account at a brokerage firm with a processor to facilitate authentication, there is still about 1% error rate, generally characterized by false rejection of legitimate users, due to the inherent imperfection of one's fingerprints. For example, if a person regularly works with abrasive chemicals, the quality of her fingerprints tends to deteriorate throughout the years. The degraded quality of the fingerprints, when faced with a security sensitive system as in most security-sensitive transactions, will certainly add to the agony of the users, thus further eroding the public's confidence toward the integrity of future systems.\nOn the other hand, if the security sensitivity is forced to be compromised to minimize false \"rejection\", then the error rate of false \"acceptance\" may increase and vice versa. Conversely, if the security sensitivity is forced to be compromised to minimize false \"acceptance,\" then the error rate of false \"rejection\" may increase. Now that a half-way decent \"match\" will allow access erroneously. This is also not something which will contribute to the public's confidence toward fingerprint-based personal identification systems. Nor will it contribute to the industry whose primary application of the fingerprint-based personal identification systems is to protect their business and financial interests.\nFurthermore, the creation of an initial file, i.e., when the account holder first sets up her account with her fingerprints at the institution's facility, may not be perfectly analyzed and stored as file data. The possibility of having less than perfect fingerprints on file makes the occurrence of false rejection/acceptance even more likely. For example, if the initial registration has a 90% accuracy, it would always be a 90% accuracy. It would still be a 90% match at best, even with a 100% accurate reading at the ATM at a later time. In other words, both ends of the overall system may contribute to the unreliability of the system.\nTherefore, it is desirable to have a personal identification system for use with fingerprint recognition front-ends to raise the percentage of accuracy, thus minimizing the security risks in connection with secured transactions.\nIt is also desirable to have a personal identification system for taking advantages of the conventional fingerprint recognition devices to provide a flexible solution in light of the various vendors of the front-end and back-end systems.\nIt is further desirable to have a fingerprint-based personal identification system which will provide an easy-to-use solution to the security issues involved in accessing the information superhighway."}
-{"text": "Thermal imaging or thermography is a recording process wherein images are generated by the use of imagewise modulated thermal energy.\nIn thermography two approaches are known:\n1. Direct thermal formation of a visible image pattern by imagewise heating of a recording material containing matter that by chemical or physical process changes colour or optical density.\n2. Formation of a visible image pattern by transfer of a coloured species from an imagewise heated donor element onto a receptor element.\nA survey of \"direct thermal\" imaging methods is given in the book \"Imaging Systems\" by Kurt I. Jacobson-Ralph E. Jacobson, The Focal Press--London and New York (1976), Chapter VII under the heading \"7.1 Thermography\". Thermography is concerned with materials which are not photosensitive, but are heat sensitive. Imagewise applied heat is sufficient to bring about a visible change in a thermosensitive imaging material.\nAccording to a direct thermal embodiment operating by physical change, a recording material is used which contains a coloured support or support coated with a coloured layer which itself is overcoated with an opaque white light reflecting layer that can fuse to a clear, transparent state wherein the coloured support is no longer masked. Physical thermographic systems operating with such kind of recording material are described on pages 136 and 137 of the above mentioned book of Kurt I. Jacobson et al.\nYet most of the \"direct\" thermographic recording materials are of the chemical type. On heating to a certain conversion temperature, an irreversible chemical reaction takes place and a coloured image is produced.\nIt has been suggested to use a thermoreducable silver source in combination with a reducing agent in a direct thermal film in order to increase the optical density in transmission of a printed image (see EP-A-537.795). Although continuous tones can be obtained by said printing method, the gradation produced by said printing method is too high resulting in only a few intermediate density levels. Fluctuations in the heat transfer from the heat source to the printing material result in a density difference of the final image. Thus, it is extremely difficult to obtain images having a uniform density profile. A direct thermal printing method moreover has the disadvantage that in the non-image places the co-reactants always remains unchanged, impairing the shelf-life and preservability.\nThermal dye transfer printing is a recording method wherein a dye-donor element is used that is provided with a dye layer wherefrom dyed portions or incorporated dye is transferred onto a contacting receiving element by the application of heat in a pattern normally controlled by electronic information signals.\nIn European Patent Application No. 94200612.3, a thermal imaging process is provided using\n(i) a donor element comprising on a support a donor layer containing a binder and a thermotransferable reducing agent capable of reducing a silver source to metallic silver upon heating and (ii) a receiving element comprising on a support a receiving layer comprising a silver source capable of being reduced by means of heat in the presence of a reducing agent, said thermal imaging process comprising the steps of\nbringing said donor layer of said donor element into face to face relationship with said receiving layer of said receiving element, PA1 image-wise heating a thus obtained assemblage by means of a thermal head, thereby causing image-wise transfer of an amount of said thermotransferable reducing agent to said receiving element in accordance with the amount of heat supplied by said thermal head and PA1 separating said donor element from said receiving element. PA1 bringing said donor layer of said donor element into face to face relationship with said receiving layer of said receiving element, PA1 image-wise heating a thus obtained assemblage preferably by means of a thermal head, thereby causing image-wise transfer of an amount of said thermotransferable reducing agents to said receiving element in accordance with the amount of heat supplied and PA1 separating said donor element from said receiving element.\nThis printing method is further referred to as `reducing agent transfer printing` or `RTP`.\nHowever, the stability of the donor element in said European Patent Application has been found to be poor. More particularly, the reducing agent tends to crystallize in the donor layer. As a result of this crystallization during storage, transfer of reducing agent is seen during printing, even on places where no heat has been applied by the thermal head. This leads to a printing fog in the final image. This problem is especially seen when a high amount of reducing agent is used in the donor layer. This high concentration is necessary to obtain high optical densities of the final printed image (above 2.0-2.5).\nMoreover, the neutral hue of the grey scale of the printed image is dependent on the choice of a specific reducing agent. It is extremely difficult to find reducing agents that yield a neutral grey-tone (e.g. for medical diagnostics). It has also been found that the reducing agent shows the disadvantage after transfer that the oxidation product of the reducing agent tends to crystallize in the receiving element, giving rise to `white dust` at the surface of the print after storage."}
-{"text": "The present invention relates to scanners, and more specifically to a carriage positioning structure for a scanner which enables the scanner to pick up the image of an object placed thereon, or turned upside down to pick up the image of an object placed below it.\nRegular scanners are commonly designed to pick up the image of flat sheets of document (2-dimensional scanning). If a sufficient depth of field is provided, a scanner can be used to pick up the image of a 3-D object. However, because regular scanners are commonly designed for 2-dimensional scanning, it is difficult to use a regular scanner to pick up the image of a 3-D object. When using a regular scanner to pick up the image of a 3-D object, the scanner may have to be turned upside down, however the image pick-up module (mechanism) cannot be stably reciprocated when the scanner is turned upside down."}
-{"text": "It is known for a rocket engine to include as a component thereof a thrust chamber having several sections joined to each other in an axial direction. Each of the sections is formed by a wall structure having an inner wall, an outer wall parallel thereto, and cooling channels formed between the inner and outer walls. The wall structure is continuous in the circumferential direction of the section.\nThe inner wall of each of two joined (to be joined) sections projects further (longer) in the extension direction of the wall structure than does the outer wall. The projecting end portion of the inner wall of one section is joined to the adjacent projecting end portion of the inner wall of the other section by a weld joint. In this way, a substantially continuous inner wall is achieved.\nThereafter, a ring-shaped element is radially arranged outside the weld joint, and said element is joined to the end portions of the adjacent outer walls. In this way, the cooling channels of one of the sections can communicate with the cooling channels of the adjacent section.\nEven though the above described rocket engine component performs satisfactorily, there is still a desire to increase the life of the component so that it can be used for an increased number of engine cycles."}
-{"text": "Insect pests are a serious problem in agriculture. They destroy millions of acres of staple crops such as corn, soybeans, peas, and cotton. Yearly, these pests cause over $100 billion dollars in crop damage in the U.S. alone. In an ongoing seasonal battle, farmers must apply billions of gallons of synthetic pesticides to combat these pests. However, synthetic pesticides pose many problems. They are expensive, costing U.S. farmers almost $8 billion dollars per year. They force the emergence of insecticide-resistant pests, and they can harm the environment.\nOther approaches to pest control have been tried. In some cases, crop growers have introduced \u201cnatural predators\u201d of the species sought to be controlled such as non-native insects, fungi, and bacteria like Bacillus thuringiensis. Alternatively, crop growers have introduced large colonies of sterile insect pests in the hope that mating between the sterilized insects and fecund wild insects would decrease the insect population. Unfortunately, success has been equivocal and the expense considerable. For example, as a practical matter, introduced species rarely remain on the treated land\u2014spreading to other areas as an unintended consequence. Predator insects migrate, and fungi or bacteria wash off of plants into streams and rivers. Consequently, crop growers need more practical and effective solutions.\nOne relatively recent solution has been to genetically engineer crops to express plant lipases that have insecticidal properties. Until now, such insecticidal lipases have only been described in certain plants, such as patatin from the potato (U.S. Pat. No. 5,743,477) and pentin from the oil bean tree (U.S. Pat. Nos. 6,057,491 and 6,339,144). However, plant-derived lipases have the inherent disadvantage of having induced natural selection pressure in insects feeding on these plants in the wild. Thus, alternative lipases are needed for insect resistance management. The present invention is useful for avoiding the inherent disadvantage of pre-existing natural selection pressure, while conferring numerous other advantages such as low cost relative to repeated-application pesticides and effective insecticidal properties."}
-{"text": "The present invention generally relates to the deglycerolization of blood, and more particularly, to a system which controls the deglycerolization of blood by monitoring the segregation of erythrocytes by size.\nThe Armed Services Blood Program Office (ASBPO) has established a policy of maintaining pre-positioned stockpiles of frozen red blood cells, and utilizing these stockpiles in times of conflict for U.S. combat casualties. In order to implement this policy, glycerol is allowed to be absorbed by red blood cells, which then are frozen and stored. The glycerol prevents damage to the erythrocytes. Presently, the only method approved by the Food and Drug Administration (FDA) for processing thawed-frozen red blood cells uses an open, nonsterile wash system that is manually monitored and operated. This system generally requires about 11/2 to 2 hours to thaw and deglycerolize red blood cells from a cryogenic state. Because this system is not sterile, the FDA mandates that thawed-frozen red blood cells processed this way must be transfused within 24 hours or discarded. However, the time restrictions and requirement to discard the blood are not compatible with the logistics of the ASBPO policy. Therefore, a need exist for a sterile, automated method for monitoring and controlling the deglycerolization of thawed red blood cells in a more timely manner compared to the processing time of the standard method."}
-{"text": "1. Technical Field\nThe present invention relates, generally, to a control system for maintaining optimum efficiency of a backlight and, more particularly in a preferred embodiment, to a closed loop temperature controller for adjusting the temperature within a fluorescent lamp to thereby optimize lamp arc drive for a given predetermined brightness set point.\n2. Background Art and Technical Problems\nScreen displays which employ fluorescent lamp backlights are used extensively in commercial, military, and consumer electronic applications. For example, such backlights are commonly used in desktop computers, laptop computers, screen displays for industrial equipment, and in connection with \"heads up\" or other screen displays in the cockpits in both commercial and military aircraft.\nConventional fluorescent lamps are commonly employed in backlit Liquid Crystal Display (LCD) applications. In a typical LCD, alphanumeric characters and other graphical images are produced on the viewing screen by selectively energizing or de-energizing preselected pixels in a two dimensional matrix to display the information. In a normally black screen display, predetermined pixels are illuminated to display the data or information as illuminated characters on a black (or other dark shade) background. In a normally white display, on the other hand, the desired data and/or information corresponds to the non-illuminated pixels, such that the information appears as black (or other dark color) images on a white (or other light color) background. In either case, a bright, consistent \"background\" light is necessary to achieve desirable contrast on the flat screen display. Indeed, in certain applications (e.g., military avionics), the high contrast provided by a bright backlight is essential to proper operation of the display.\nIt is also desirable to obtain a desired brightness while minimizing power consumption. This is particularly important in portable electronics, for example laptop computers and the like, where battery life is an important product feature.\nPresently known systems for controlling the brightness of a fluorescent backlight lamp typically involve a control system for supplying lamp arc drive to the backlight, to thereby excite the gas atoms within the sealed lamp enclosure to create visible light. The amount of visible light emitted by the lamp is sensed, for example by a photodiode, and a feedback signal indicative of the brightness output of the lamp is fed back to a control circuit. This feedback signal (indicative of actual brightness) is compared to an input signal representative of a desired brightness level, and presently known control systems drive the difference between this actual signal and the desired signal to a minimum. Under this control regime, if the actual brightness is less than the desired brightness, the controller increases the lamp arc drive applied to the lamp until the actual brightness equals the desired brightness. If, on the other hand, the actual brightness is greater than the desired brightness, the controller circuit reduces the magnitude of the lamp arc drive applied to the lamp until the actual brightness emitted from the lamp again equals the desired brightness level for the lamp. Presently known prior art brightness control systems typically employ a \"cold spot\" at a predetermined point on the lamp which functions to keep a certain amount of the gas (typically mercury) within the lamp in a condensed state. Such \"cold spot\" systems employ the well known principle that maintaining the temperature of the cold spot in a specified range allows for very efficient operation of the lamp. Presently known systems, however, often require expensive components to maintain the cold spot, and do not adequately compensate for drifting or degradation over time of some of the parameters which influence the efficiency of the lamp.\nA fluorescent lamp control system is thus needed which overcomes the shortcomings of the prior art."}
-{"text": "The present invention relates to magnetic resonance imaging (MRI) systems, and particularly to the radio frequency (RF) coils used in such systems.\nMagnetic resonance imaging (MRI) utilizes hydrogen nuclear spins of the water molecules in the human body or other tissue, which are polarized by a strong, uniform, static magnetic field generated by a magnet (referred to as B0\u2014 the main magnetic field in MRI physics). The magnetically polarized nuclear spins generate magnetic moments in the human body. The magnetic moments point in the direction of the main magnetic field in a steady state, and produce no useful information if they are not disturbed by any excitation.\nThe generation of nuclear magnetic resonance (NMR) signal for MRI data acquisition is achieved by exciting the magnetic moments with a uniform radio frequency (RF) magnetic field (referred to as the B1 field or the excitation field). The B1 field is produced in the imaging region of interest by an RF transmit coil which is driven by a computer-controlled RF transmitter with a power amplifier. During the excitation, the nuclear spin system absorbs magnetic energy, and its magnetic moments precess around the direction of the main magnetic field. After the excitation, the precessing magnetic moments will go through a process of free induction decay, emitting their absorbed energy and then returning to the steady state. During the free induction decay, NMR signals are detected by the use of a receive RF coil, which is placed in the vicinity of the excited volume of the human body. The NMR signal is an induced electrical motive force (voltage), or current, in the receive RF coil that has been induced by the flux change over some time period due to the relaxation of precessing magnetic moments in the human tissue. This signal provides the contrast information of the image. The receive RF coil can be either the transmit coil itself, or an independent receive-only RF coil. The NMR signal is used for producing magnetic resonance images by using additional pulsed magnetic gradient fields, which are generated by gradient coils integrated inside the main magnet system. The gradient fields are used to spatially encode the signals and selectively excite a specific volume of the human body. There are usually three sets of gradient coils in a standard MRI system, which generate magnetic fields in the same direction of the main magnetic field, varying linearly in the imaging volume.\nIn MRI, it is desirable for the excitation and reception to be spatially uniform in the imaging volume for better image uniformity. In a standard MRI system, the best excitation field homogeneity is usually obtained by using a whole-body volume RF coil for transmission. The whole-body transmit coil is the largest RF coil in the system. A large coil, however, produces lower signal-to-noise ratio (S/N) if it is also used for reception, mainly because of its greater distance from the signal-generating tissues being imaged. Since a high signal-to-noise ratio is the most desirable factor in MRI, special-purpose coils are used for reception to enhance the S/N ratio from the volume of interest.\nIn practice, a well-designed specialty RF coil should have the following functional properties: high S/N ratio, good uniformity, high unloaded quality factor (O) of the resonance circuit, and high ratio of the unloaded to loaded Q factors. In addition, the coil device must be mechanically designed to facilitate patient handling and comfort, and to provide a protective barrier between the patient and the RF electronics. Another way to increase the S/N is by quadrature reception. In this method, NMR signals are detected in two orthogonal directions, which are in the transverse plane or perpendicular to the main magnetic field. The two signals are detected by two independent individual coils which cover the same volume of interest. With quadrature reception, the S/N can be increased by up to \u221a2 over that of the individual linear coils.\nTo cover a large field-of-view, while maintaining the S/N characteristic of a small and conformal coil, a linear surface coil array technique was created to image the entire human spines (U.S. Pat. No. 4,825,162). Subsequently, other linear surface array coils were used for C.L. spine imaging, such as the technique described in U.S. Pat. No. 5,198,768. These two devices consist of an array of planar linear surface coil-elements. These coil systems do not work well for imaging deep tissues, such as the blood vessels in the lower abdomen, due to sensitivity drop-off away from the coil surface.\nTo image the lower extremities, quadrature phased array coils have been utilized such as described in U.S. Pat. Nos. 5,430,378 and 5,548,218. The first quadrature phased array coil, images the lower extremities by using two orthogonal linear coil arrays: six planar loop coil elements placed in the horizontal plane and underneath the patient and six planar loop coil elements placed in the vertical plane and in between the legs. Each linear coil array functions in a similar way as described in U.S. Pat. No. 4,825,162 (Roemer). The second quadrature phased array coil (Lu) was designed to image the blood vessels from the pelvis down. This device also consists of two orthogonal linear coil arrays extending in the head-to-toe direction: a planar array of loop coil elements laterally centrally located on top of the second array of butterfly coil elements. The loop coils are placed immediately underneath the patient and the butterfly coils are wrapped around the patient. Again, each linear coil array functions in a similar way as described in U.S. Pat. No. 4,825,162.\nIn MRI, gradient coils are routinely used to give phase-encoding information to a sample to be imaged. To obtain an image, it is required that all the data points in a so-called \u201ck-space\u201d (i.e., frequency space) must be collected. Recently, there have been developments where some of the data points in k-space are intentionally skipped and at the same time use the intrinsic sensitivity information of RF receive coils as the phase-encoding information for those skipped data points. This action takes place simultaneously, and thus is referred to as partially parallel imaging or partially parallel acquisition (PPA). By collecting multiple data points simultaneously, it requires less time to acquire the same amount of data, when compared with the conventional gradient-only phase-encoding approach. The time savings can be used to reduce total imaging time, in particular, for the applications in which cardiac or respiratory motions in tissues being imaged become concerns, or to collect more data to achieve better resolution or S/N. SiMultaneous Acquisition of Spatial Harmonics, SMASH, (U.S. Pat. No. 5,910,728 and \u201cSimultaneous Acquisition of Spatial Harmonics (SMASH): Fast Imaging with Radiofrequency Coil Arrays,\u201d Daniel K. Sodickson and Warren J. Manning, Magnetic Resonance in Medicine 38:591\u2013603 (1997), both incorporated herein by reference) and \u201cSENSE: Sensitivity Encoding for Fast MRI,\u201d Klaas P. Pruessmann, et al., Magnetic Resonance in Medicine 42:952\u2013962 (1999, also incorporated by reference, are basically two methods of PPA. SMASH takes advantage of the parallel imaging by skipping phase encode lines that yield decreasing the Field-of-View (FOV) in the phase-encoding direction and uses coils (e.g., coil arrays) together with reconstruction techniques to fill in the missing data points in k-space. SENSE, on the other hand, is a technique that utilizes a reduced FOV in the read direction, resulting an aliased image that is then unfolded in x-space (i.e., real space), while using the RF coil sensitivity information, to obtain a true corresponding image. Here, we make use of phase difference between signals from multiple coils to skip phase encoding steps. By skipping some of the phase encoding steps, one can achieve speeding up imaging process by a reduction factor R. Theoretically speaking, the factor R should equal the number of independent coils/arrays. In the SENSE approach, the SIR is defined as:                SNRSENSE=SNRFULL/{g\u221aR}where SNRFULL is the S/N achievable when all the phase encoding steps are collected by traditional gradient phase encoding scheme. SNRSENSE is optimized when the geometry factor g equals 1. To obtain g of 1, traditional decoupling techniques such as overlapping nearest neighbor elements to null the mutual inductance between them shall not apply, as have been reported by others.        \nSENSE and SMASH or a hybrid approach of both demand a new type of design requirements in RF coil design. In SMASH, the primary criterion for the array is that it be capable of generating sinusoids whose wavelengths are on the order of the FOV. This is how the target FOV along the phase encoding direction for the array is determined. Conventional array designs can incorporate element and array dimensions that will give optimal S/N for the object of interest. In addition, users of conventional arrays are free to choose practically any FOV, as long as severe aliasing artifacts are not a problem. In contrast, when using SMASH, the size of the array determines the approximate range of FOVs that can be used in the imaging experiment. This then determines the approximate element dimensions, assuming complete coverage of the FOV is desired, as in most cases. In SENSE, the method is based upon the fact that the sensitivity of a RF receiver coil generally has a phase-encoding effect complementary to those achieved by linear field gradients. For SENSE imaging, the elements of a coil array should be smaller than for common phased-array imaging, resulting in a trade-off between basic noise and geometry factor, and adjacent coil elements should not overlap for a net gain in S/N due to the improved geometry factor when using SENSE.\nFor PPA applications, different types of RF coils or arrays have been used so far. However, most of them are based upon \u201ctraditional\u201d RF coil design requirements, thus remain within the conventional coil design scheme. It has been reported, however, that since the phase information of B1 of a receive coil is very important when SENSE applications are demanded, for example, new coil design techniques such as non-overlapping adjacent coil elements may be necessary for better definition of the individual phase information associated with each RF coil used in an array, unlike traditional design scheme where two adjacent coils elements are overlapped to null the mutual inductance between the elements (U.S. Pat. No. 4,825,162). Without overlap, the coupling may be increased, but there is a net gain in S/N due to the improved geometry factor when using SENSE. As stated in the above, the use of smaller coil-elements than those for conventional imaging results in a trade-off between basic noise and geometry factor."}
-{"text": "There are many industrial processes and applications that require a continuous supply of a process fluid used for cooling purposes. The devices used to furnish this process fluid supply are ordinarily called \"water chillers\", even though the process fluid is most usually a mixture of water and other liquid, usually ethylene glycol, and that term is used in this specification to refer to controlled chillers using any appropriate coolant.\nA water chiller must maintain the temperature of the process fluid essentially constant within a very limited tolerance range; in some instances the critical parameter may be selected as the return temperature from the process apparatus, but more frequently the supply temperature should be selected as the basis of control. The water chiller is often a self-contained unit that can be transported from one location to another, literally constituting a heat transfer system on wheels. Water chillers of this kind are most frequently employed in processes involving heat transfer where the coolant must be maintained constant at some temperature between -30.degree. F and +60.degree. F.\nFor the most part, the controls for water chillers have constituted conventional thermostats, freezestats, and pressure-actuated switches connected directly in the electrical actuating circuits of the pumps, compressors, and other components of the water chilling apparatus. A further control that is often used is a float switch, connected to a coolant reservoir, employed to interrupt operation whenever the water supply is inadequate. Another control frequently used is a hot-gas bypass valve, used to control cooling capacity during intermittent or partial load conditions.\nFor many applications, these direct acting controls are adequate. However, they do not afford the precision control essential to some sensitive manufacturing and other industrial operations. Furthermore, controls of this nature do not provide comprehensive protection for the water chiller and the process apparatus it serves relative to possible malfunctions in the chiller itself or in the process apparatus. Moreover, the inherent inertia of these conventional controls makes it difficult to control chiller operation on the basis of coolant supply temperature; it is usually necessary to use the coolant return temperature as the control parameter."}
-{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device of a chip-on-chip structure that has a semiconductor chip bonded with other semiconductor chips thereon, and to a method for manufacturing same.\n2. Description of Related Art\nThere are semiconductor devices increased in integration degree, including system-on-chips (SOCs) and multi-chip modules (MCMs). In the system-on-chip, the functions conventionally realized on a plurality of ICs are integrated on one semiconductor chip. Meanwhile, the multi-chip module is structured by a plurality of semiconductor chips arranged with density on a wiring board of glass-epoxy or the like. Each of them has multiple functions as one semiconductor device and can be size-reduced as compared to a combination of a plurality of semiconductor devices realizing the equivalent functions. Meanwhile, this reduces the wiring length in the overall, enabling high-speed transmission of signals.\nHowever, the system-on-chip is complicated in manufacture process, requiring a huge amount of capital investment and hence high manufacture cost. Meanwhile, the multi-chip-module has a plurality of semiconductor chips mutually arranged side by side on a wiring board. Because these semiconductor chips are connected by wirings, the size is greater as compared to the system-on-chip and hence integration degree lowers.\nOn the other hand, there is a chip-on-chip structure as one form of an integration increased semiconductor device. The chip-on-chip-structured semiconductor device has a structure having a plurality of semiconductor chips that are oppositely placed for mutual connection. Such a semiconductor device is not configured with an integration of the functions conventionally realized on a plurality of ICs (semiconductor chips) as in the system-on-chip. Consequently, the manufacture process is not so complicated as that of the system-on-chip. Therefore, there is a merit to reduce the manufacturing cost.\nThe chip-on-chip-structured semiconductor devices include those having a plurality of small semiconductor chips (child chips) arranged side by side on one large semiconductor chip (parent chip). Such a semiconductor device, at a glance, is similar in structure to the multi-chip module having a plurality of semiconductor chips arranged side by side on a wiring board.\nIn the chip-on-chip-structured semiconductor device, however, the parent chip not only serves as a wiring board to connect between the child chips but also acts itself as a semiconductor chip having functional elements. This results in higher integration degree. On the other hand, the wiring formed on the parent chip, made by a semiconductor process, is by far finer than the wiring of a wiring board of a multi-chip-module. This allows the functional elements of the semiconductor chips (parent and child chips) to be connected at a short length of wiring, whereby signal transmission speed can be increased higher as compared to that of the multi-chip module.\nThe chip-on-chip-structured semiconductor devices include those that a child chip has, further, one or a plurality of child chips vertically superposed thereon. Namely, such a semiconductor device has a structure connected, on a parent chip, with one or a plurality of chip blocks each having one or a plurality of superposed semiconductor chips. This structure realizes a highly integrated semiconductor device.\nIn such a semiconductor device, however, the wiring between arbitrary two of the semiconductor chips is provided necessarily through a wiring plane (usually, active surface) of the parent chip, thereby increasing the mean wiring length. Namely, in the case the semiconductor chip is in a position close to the top end of the chip block (position distant from the parent chip), there is an increased length of wiring between that semiconductor chip and the parent chip. Consequently, signals could not be transmitted at sufficient high speed. Meanwhile, even if the wiring length is tried to decrease throughout the entire semiconductor device, there is encountered a low freedom in design, e.g. restriction in semiconductor chip arrangement."}
-{"text": "1. Field of the Invention\nThe present invention relates to a club-head for a golf club and, more particularly, to an improvement of a club-head having a hitting portion which includes a main body made of a fiber-reinforced plastic and a metal sole member integrally fixed to the main body along the underside of the main body.\n2. Description of the Related Arts\nRecently, a club-head has been used in which a hitting portion provided at the heel thereof with a hosel or neck portion for connecting a shaft consists of a main body made of a fiber-reinforced plastic for hitting a ball at the front surface thereof and a sole member made of a metal is integrally joined to the main body along the underside thereof. The club-head having such a construction has a drawback in that the main body is apt to be exfoliated from the sole member upon impact with a golf ball.\nU.S. patent application No. 840,795 filed by the present applicant on Mar. 18, 1986, discloses one kind of the above-mentioned club-head wherein the back weight plate made of a heavy metal is attached to or embedded in the rear surface of the main body above the sole member. It has known that the club-head having such a construction can increase a flight distance of a golf ball hit by the head because the heavy back weight plate can most efficiently serve the main body to effect a repulsion action on the golf ball. But, the club-head having such a construction also has the drawback in that the main body is apt to be exfoliated from the sole member upon impact with a golf ball.\nU.S. patent application No. 854,626 filed by the present applicant on Apr. 22, 1986, discloses another kind of the above-mentioned club-head wherein the weight member is embedded in the main body and connected to the sole member through one or more connecting members. In this construction, when the heavy weight member is arranged between the front and rear surfaces of the hitting portion of the club-head, the thickness of the main body between the front surface of the hitting portion of the club-head and the weight member is decreased, and thus, the repulsion action of the head against the golf ball and depth of the center of gravity of the head are decreased. On the other hand, when the heavy weight member is located at the rearmost position of the hitting portion of the club-head, the head has a drawback in that the main body is apt to be exfoliated from the sole member."}
-{"text": "Today's modern aircraft use hermetically sealed fire extinguishers that are opened, or activated, by direct explosive impingement energy. With reference to Prior Art FIG. 1, the device which provides the explosive energy is typically called a cartridge 120, or squib. The impingement energy is focused on a dome-shaped hermetic burst disc 110 such that the burst disc will rupture as a result of the impingement. The burst disc material used is typically fabricated from corrosion resistant steel.\nTypically, the cartridge 120 is retained in a discharge head 130 in such a manner that it directly faces the burst disc 110 assembly. The discharge head 130 is attached to the outlet of the fire extinguisher and is typically used to direct the flow of extinguishing agent to an aircraft interface, such as plumbing or tubing, which directs the agent to the desired location. A filter screen 150 is located within the discharge head to catch any large burst disc fragments created as a result of the explosive impingement energy."}
-{"text": "Overall implementation of a keypad for a mobile device, i.e., the final package, including domes, backlight and keyset, should provide the required visual functionality as well as sufficient tactile feedback for the user when pressing an individual key. In view of thin keypads, the tactile feedback can be obtained using elastic light guide, dome pin lead throughs (holes above domes), or dome pins integrated with a plastic light guide (flexible but not elastic). Of these, sufficient optical quality can be obtained using either the first or the last approach.\nSpatially selective out-coupling of light from a planar light guide plate for keypad illumination is an important desired feature. It means that simultaneous control of out-coupling efficiency and spectra of the out-coupled light needs to be provided to enable modification of the visual appearance (intensity, colored/white) of the individual keys. In addition, a tactile feedback when pressing an individual key should not be compromised by the illumination system.\nVarious types of out-coupling means based on surface structuring of the planar light guide exist. Spatial modulation of the out-coupling can be obtained by limiting/modulating the spatial structuring to/within separate regions (key areas) on the light guide.\u2019\nMoreover, colored appearance can be typically obtained using white light out-coupling and \u2018painted\u2019 keysets on a keymat. In cost sensitive applications, such as mobile phones, the utilized \u2018paints\u2019 typically have large variety of optical quality depending on the utilized manufacturer and can cause significant efficiency loss due to poor optical characteristics.\nVarious difficulties appear in keypad utilizing light guides for illumination and usually requiring several different manufacturing steps in keypad and light guide fabrication. Manufacturing of keypad's dome pins is usually realized by injection molding. The mechanical key function is realized by dome pins through holes in the illumination light guide. This can lead to energy losses in the illumination."}
-{"text": "There is increasing interest in IP telephony to help lower costs and enable new services. Many enterprises and call centers are adopting IP telephony over their converged IP infrastructure and many multi-site corporations are using Voice over IP (VoIP) for their intra- and inter-site communication.\nWith the use of VoIP for mission-critical business applications, it is important to evaluate and improve the reliability and quality of VoIP calls. Ideally, a VoIP call should be as reliable as a traditional circuit-switched phone call. However, when monitoring and evaluating the quality of a VoIP call (especially over wide area links), it is necessary to deal with the inherent packet losses, delays, and jitter associated with IP networks, which are not encountered in traditional circuit-switched networks. Even though IP networks are largely self-healing for network faults, and many enterprise networks are engineered to have redundant links or paths between sites, today's IP networks are not engineered to react to performance degradations at the timescales needed for voice. For example, recent studies show that while there is acceptable performance within some service provider networks, many backbone paths still have poor VoIP performance and network faults cause problems.\nTo provide a robust VoIP infrastructure, it is important to rapidly detect performance degradations and faults. This detection is complicated by several factors. On a per-connection basis, for example, there are natural silence periods in VoIP calls during which packets are not transmitted by a source (e.g., when a participant in a call is listening rather than speaking). Consequently, while monitoring a VoIP call (e.g., on the receiving side) it is necessary to distinguish between gaps that occur due to natural speech silences and perhaps speech compression, and the gaps that occur due to packet loss, delay, and jitter in the IP network. Furthermore, although detecting problems can certainly help to alert a network manager, it would be particularly useful if the network could react to a detected problem and route around it.\nAccordingly, there is a need for techniques for performing rapid fault detection and recovery in communication networks such as IP telephony networks, particularly those that provide VoIP applications."}
-{"text": "Advances in electronics, batteries and packaging technology have led to continued increases in the number of mobile computing devices in service. These mobile computing devices often have an associated docking station that permits ready access to printers, modems, networks, and peripherals that are more humanly comfortable, amongst other connections and attributes. Though beneficial for their intended purpose, these docking stations are to some extent disadvantageous as discussed below.\nMobile computing devices tend to be small to facilitate efficient mobility. Docking stations tend to be configured for desktop use in which a mobile computing device is docked at the station for data transmission or for running the desktop environment with the resources of the mobile computing device. When docked at a desk top station, the screen of a mobile computing device is small, located at a relatively far distance from an operator and positioned such that an operator often has to place his or her back, neck or head in an uncomfortable position to adequately see the screen.\nErgonomic studies of computer-human operator interfaces have determined that a preferred monitor or screen height is that at which the top of the screen is even with the horizontal line of sight of the operator. The preferred angle of the screen is at 90 degrees to the line of sight.\nOne attempt to alleviate the undesirably low screen height of docked mobile computing devices includes providing a regular desk top monitor on a stand above the docking station. Though this may alleviate some of the problems of low-level mobile computing device screen height, the additional monitor and stand are often undesirably expensive and consumptive of space. This arrangement may also provide insufficient adjustment of screen viewing angle."}
-{"text": "Field of the Invention\nThe present invention relates to a conductive material and a substrate having a conductive film formed thereon by using the conductive material.\nDescription of the Related Art\nA polymer having a conjugated double bond (i.e. \u03c0-conjugated polymer) does not show a conductivity by itself; however, if an appropriate anionic molecule is doped therein, it can express a conductivity, thereby giving a conductive polymer material (i.e. conductive polymer composition). As to the \u03c0-conjugated polymer, polyacetylene, (hetero) aromatic polymers such as polythiophene, polyselenophene, polytellurophene, polypyrrole, and polyaniline; a mixture thereof, etc., are used; and as to the anionic molecule (dopant), an anion of sulfonic acid type is most commonly used. This is because a sulfonic acid, which is a strong acid, can efficiently interact with the aforementioned \u03c0-conjugated polymers.\nAs to the anionic dopant of sulfonic acid type, sulfonic acid polymers such as polyvinyl sulfonic acid and polystyrene sulfonic acid (PSS) are widely used (Patent Document 1). The sulfonic acid polymer includes a vinylperfluoroalkyl ether sulfonic acid typified by Nafion (registered trademark), which is used for a fuel cell.\nPolystyrene sulfonic acid (PSS) has a sulfonic acid as a repeated monomer unit in the polymer main chain, so that it has a high doping effect to the \u03c0-conjugated polymer, and also can enhance water dispersibility of the \u03c0-conjugated polymer after being doped. This is because the hydrophilicity is kept due to the sulfo groups excessively present in PSS, and the dispersibility into water is therefore enhanced dramatically.\nPolythiophene having PSS as a dopant exhibits high conductivity and can be handled as an aqueous dispersion, so that it is expected to be used as a coating-type conductive film material in place of ITO (indium-tin oxide). As mentioned above, however, PSS is a water-soluble resin, and is hardly soluble in an organic solvent. Accordingly, the polythiophene having PSS as a dopant also has a high hydrophilicity, but a low affinity to an organic solvent and an organic substrate, and thus, it is difficult to disperse it into an organic solvent or to form a film onto an organic substrate.\nBesides, when the polythiophene having PSS as a dopant is used in, for example, a conductive film for an organic EL lighting, a large quantity of water tends to remain in the conductive film and the conductive film thus formed tends to absorb moisture from an outside atmosphere since the polythiophene having PSS as a dopant has an extremely high hydrophilicity as mentioned above. As a result, the problems arise that the luminous body of the organic EL chemically changes, thereby the light emitting capability is deteriorated, and that water agglomerates over time and defects are caused, which results in shortening of the lifetime of the whole organic EL device. Furthermore, there arise other problems in the polythiophene having PSS as a dopant that particles in the aqueous dispersion becomes large, the film surface becomes rough after the film formation, and a non-light emitting region, called dark spot, is caused when used for the organic EL lighting.\nIn addition, since the polythiophene having PSS as a dopant has an absorption at a wavelength of about 500 nm in the blue region, in the case that this material is used as a film coating a transparent substrate such as a transparent electrode, there arises another problem that when the conductivity required for the device to function is made up by the solid concentration or the thickness of the film, transmittance of the film is affected.\nPatent Document 2 discloses a conductive polymer composition composed of a conductive polymer which contains a \u03c0-conjugated polymer formed of a repeating unit selected from thiophene, selenophene, tellurophene, pyrrole, aniline, and a polycyclic aromatic compound, and a fluorinated acid polymer which can be wetted by an organic solvent and 50% or more of which is neutralized by a cation; and it is shown that an aqueous dispersion of the conductive polymer can be obtained by combining water, a precursor monomer of the \u03c0-conjugated polymer, the fluorinated acid polymer, and an oxidant, in any order.\nHowever, in such a conventional conductive polymer, particles are agglomerated in the dispersion immediately after synthesis. Also, if an organic solvent served as a conductive enhancer is added thereto to give a coating material, the agglomeration is further facilitated, so that the filterability thereof is deteriorated. If the coating material is applied by spin coating without filtration, a flat film cannot be obtained due to the effect of the particle agglomeration; and as a result, the problem of coating defect is caused.\nMoreover, development has been promoted in a flexible device. As a transparent conductive film for the current hard devices, ITO is widely used. ITO is, however, a crystalline film, and therefore there arises cracks in case of bending. Accordingly, it is a pressing need to develop a flexible transparent conductive film substituting for ITO. Polythiophene having PSS as a dopant forms a flexible film with high transparency, but involves a problem of low conductivity compared to ITO in addition to the aforementioned problem of dark spot.\nPatent Document 3 discloses a transparent conductive film using silver nanowires. The transparent conductive film using silver nanowires is one of a candidate for a conductive film for a flexible device, since it has high conductivity and transparency. The film using silver nanowire, however, conducts electricity only through the wire part, and therefore causes a problem that the light emission occurs at the wire parts only, not the whole surface when it is applied to an organic EL lighting."}
-{"text": "In a typical ink jet recording or printing system, ink droplets are ejected from a nozzle at high speed towards a recording element or medium to produce an image on the medium. The ink droplets, or recording liquid, generally comprise a recording agent, such as a dye or pigment, and a large amount of solvent. The solvent, or carrier liquid, typically is made up of water, an organic material such as a monohydric alcohol, a polyhydric alcohol or mixtures thereof.\nA continuing problem with ink jet printers is the accumulation of ink on ink jet nozzle plates, particularly around the orifice from which ink drops are ejected. The result of ink drops accumulating around the orifice is that it becomes wettable causing ink drops to be misdirected, degrading the quality of the printed image.\nTo limit or prevent the spreading of ink from the orifice to the nozzle plate, it is common practice to coat the ink jet nozzle plate with an anti-wetting layer. Examples of anti-wetting layers are coatings of hydrophobic polymer materials such as Teflon.RTM. and polyimide-siloxane, or a monomolecular layer of a material that chemically binds to the nozzle plate, e.g., alkyl thiols, alkyl trichlorosilanes and partially fluorinated alkyl silanes.\nInk jet nozzle plates are also contaminated by ink drops that land on the nozzle plate. These \"satellite\" ink drops are created as a by-product of the drop separation process of the primary ink drop that is used to print. Another source of contaminating ink are tiny ink drops that are created when the primary ink drop impacts recording material. Ink drops accumulating on the nozzle plate can also potentially attract contaminants such as paper fibers which cause the nozzles to become blocked. Partially or completely blocked nozzles can lead to missing or misdirected drops on the print media, either of which degrades the quality of the print.\nIn order to solve this problem, the nozzle plates are periodically wiped clean. Several wiping methods are known including wet wiping techniques utilizing inks as a cleaning solvent. While inks and ink solvents used to dilute inks may be used as a cleaning liquid, they are not optimized for this purpose. Inks may contain additives such as, for example, ethylene glycol, diethylene glycol, and diethylene glycol monobutyl ether which may be environmentally undesirable when released during cleaning in unventilated areas such as a home or an office.\nFurther, inks often contain various materials which may leave an undesirable residue on the ink jet print head nozzle plate. Thus while wiping removes ink drops from the nozzle plate, the hydrophobic anti-wetting coating on the nozzle plate may be severely contaminated and compromised by ink residue. The ink-fouled coating is therefore unable to prevent the spreading of ink from orifices.\nIt has also been discovered that hydrophobic coatings on an ink jet print head nozzle plate are susceptible to fouling by certain ink jet inks, such as those containing copper phthalocyanine dyes. The fouling of the nozzle plate by the ink can lead to excessive spreading by ink on to the nozzle plate during normal use, further aggravating drop placement problems. Another disadvantage in using inks as a cleaning solution is that they are expensive.\nThere remains a need for a simple, economical ink jet nozzle plate cleaning solution that will help maintain the anti-wetting character of ink jet nozzle plates so that an ink jet print head will consistently deliver accurate and reproducible drops of ink to a receiver resulting in photographic quality images."}
-{"text": "The present invention relates to a semiconductor device and a method of manufacturing the same, more specifically, a semiconductor device comprising a nonvolatile memory of the stacked gate structure and a transistor of the single-layer gate structure, and a method of manufacturing the semiconductor device.\nThe logic semiconductor device combined with a nonvolatile semiconductor memory forms product fields, as of CPLD (Complex Programmable Logic Device) and FPGA (Field Programmable Gate Array), and because of their characteristic of programmability, so far have formed large markets because of their characteristic, programmability.\nThe logic semiconductor device combined with a nonvolatile memory has, in addition to flash memory cells, a high-voltage transistors for controlling the flash memory and low-voltage transistors of high-performance logic circuit integrated on the same semiconductor chip. The flash memory cells have gate electrode of the stacked structure of a floating gate and a control gate laid one on the other which is different from the single-layer structure of the high-voltage transistors and the low-voltage transistors. Accordingly, the process of manufacturing the logic semiconductor device combined with the nonvolatile memory requires the process specialized in forming the nonvolatile memory transistors of the stacked gate structure without changing characteristics of the peripheral circuits, especially the logic transistors.\nIn the usual combined process, the floating gates of the nonvolatile memory transistors are formed of the first-level conductive film (the first conductive film), and the control gates of the nonvolatile memory transistors and the gate electrodes of the peripheral transistors are formed of the second-level conductive film (the second conductive film). Then, the peripheral transistors are formed after the nonvolatile memory transistors have been formed, so as to prevent the process of manufacturing the nonvolatile memory transistors from influencing characteristics of the logic transistors. In terms of the process of forming the gate electrodes, after the second conductive film in the memory cell region have been patterned to form the control gates, the second conductive film in the peripheral circuit region is patterned to form the gate electrodes of the peripheral transistors.\nThe related arts are described in, e.g., Reference 1 (Japanese published unexamined patent application No. Hei 10-209390).\nHowever, the inventors of the present application have examined the process of fabricating the logic semiconductor device combined with the nonvolatile memory and found that the process causes the disadvantage that the second conductive film in the memory cell region is etched when the second conductive film is patterned to form the gate electrodes of the peripheral transistors."}
-{"text": "This invention relates to a method for laminating a transparent safety panel to the screen-bearing viewing window of a CRT (cathode-ray tube) and particularly, but not exclusively, to a large CRT, and to the product of that method. By \"large CRT\" is meant a television picture tube or an information display tube having a viewing window bearing a viewing screen of at least a 30V size.\nIn one prior type of CRT, a glass safety panel is laminated to the viewing window of the CRT in order to reduce the danger of implosion and also, should the CRT implode, to reduce the danger of injury to people and things nearby. Suitable structures and methods for laminating CRTs smaller than 26V screen dimension have been described previously, for example, in U.S. Pat. No. 4,204,231 issued May 20, 1980 to M. M. Permenter.\nIn one prior laminating method, a safety panel is positioned in a desired spaced-apart relationship with a CRT window, and then a strip of flexible pressure-sensitive tape is wrapped around the edges of the CRT and panel to define a closed cell between the window and the panel. Thereafter, the cell is filled with a catalyzed liquid resin and allowed to cure to a clear transparent solid which adheres to the facing surfaces of the panel and the window. During the filling step, vent holes are punctured in the tape to allow air in the cell to escape. After the cell is filled with resin, the holes are taped shut to prevent both the leakage of resin and the formation of bubbles in the cell during the curing step. A foam tape with adhesive on both sides adhered to the margins of the panel and the window may replace the tape in the foregoing method.\nIn the foregoing method, the cell is filled with the viewing window positioned in a generally horizontal attitude, with the window facing downward. The tape provides a temporary hydraulic seal for the cell, and is also sufficiently strong to hold the panel temporarily in the desired downward-facing, spaced-apart relationship while the cell is being filled and the resin cures. In practice, especially when safety panels are laminated to CRT windows larger than about 25V-size, the taping step is not efficient and many temporary seals leak resin, and some seals fail to hold the safety panels in place. Also, because the windows face downward, it is difficult to determine whether gas bubbles are present in the viewable area in front of the window.\nBecause of the size and weight of a large CRT as defined above, all of these problems are aggravated and these prior methods are not practical for laminating a safety panel to the window of a large CRT. For example, a 25V-size CRT weighs about 55 pounds, while a 33V-size CRT weighs about 130 pounds and therefore cannot be handled manually in the factory. A cell, formed as described above, on a 33V-size CRT is difficult to fill with liquid resin with the window facing down because of the greater weight of the safety panel and because the greater weight of the resin causes greater leakage of resin during the filling and curing steps, especially through the venting holes in the tape. The tapes used in the prior methods to form the cell are not reliable to hold an 11-pound safety panel and about 8 pounds of liquid resin in the desired spaced apart relationship during the filling and curing steps. Sagging and wedging of the panel position and leakage and \"flow-out\" of the resin occur frequently with prior methods."}
-{"text": "1. Field of the Invention\nThis invention relates to the art of dispensing pressure sensitive labels.\n2. Brief Description of the Prior Art\nVarious U.S. Pat. Nos. 1,642,387, 2,259,358, 2,275,064, 2,516,487, 2,620,205, 3,051,353, 3,265,553, 3,501,365, 3,551,251, and 3,611,929 and British Pat. No. 1,057,126, Feb. 1, 1967 are made of record. U.S. Pat. No. 3,501,365 British Pat. No. 1,057,126 disclose U-shaped cuts in the supporting web or strip. In connection with one embodiment of U.S. Pat. No. 3,501,365 for example, the patent discloses \"that the feed holes are die cut in such a manner that the leading edge of the feed hole has been cut whereas the trailing edge has not\" been cut. The loose flap or flap portions formed by the U-shaped cut extends in a leading direction as the flap portion moves toward the peel edge or separator point. Due to the fact that the flap portion is adhered to the adhesive on the overlying label and due to the fact that the flap portions extend in the leading direction, the flap portions fold out of the plane of the supporting web at the peel edge. Also, the flap portions adhere to the overlying label so well because when the flap-forming cuts are made, the supporting web and the label are driven together by the cutting knives into intimate contact with the intervening adhesive. This increases the holding force of the flap portions to the overlying labels and consequently diminishes the tendency of the flap portions to separate from the overlying labels at the peel edge. This sometimes results in tearing at the ends of the U-shaped cut. As the supporting material at the ends of a U-shaped cut tears, the flap portion grows in area. As the flap portion grows, the enlarged flap portion is held to the overlying label or labels by an increased area of adhesive. Accordingly, the force holding the enlarged flap portion to the overlying label or labels increases as the tearing progresses until eventual rupture or breakage of the supporting web. In another embodiment of U.S. Pat. No. 3,501,365, loose internal portions formed by annular cuts move out of the plane of the supporting web at the peel edge along with the leading label as the supporting web is drawn about the peel edge. These loose internal portions cannot tear the supporting web because they are severed therefrom by the annular cut. However, in both embodiments, the flap portions move out of the plane of the supporting web at the peel edge for the same reasons."}
-{"text": "Solid wood doors have been manufactured for many years with the same assembly principle.\nIn the last few years, the insulation of panels has become a standard step in the manufacturing of doors. The insulation of panels helps eliminate the problems of condensation on the interior of the panels, a location where the panels are the thinnest.\nIn construction, the R-value is the measurement of a material's capacity to resist heat flow from one side to the other. In simple terms, R-values measure the effectiveness of insulation and a higher number represents more effective insulation. Despite the insulation of panels, the global R value of a solid wood slab door remains relatively low.\nIn order to create a solid wood slab door with a higher R value, a system for insulating solid wood is herein provided."}
-{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device, and more particularly, to a semiconductor device that can have power consumption reduced.\n2. Description of the Background Art\nSemiconductor devices incorporated in various equipments have the scale of integration increased in order to reduce the size of the device and to integrate much more powerful logic. Increase in the integration density will result in a great number of elements operating inside to raise the heat. Therefore, reducing power consumption is an important factor. For example, in a DRAM (Dynamic Random Access Memory), there is a great demand for reducing the power consumption as a result of increase in the number of elements according to increase in the storage capacity.\nA DRAM as a conventional semiconductor memory device will be described in detail hereinafter. A DRAM includes an intermediate potential generation circuit for generating precharge potential for a bit line, a timer circuit for carrying out a self refresh operation, and an internal high voltage circuit for generating high potential to be provided to a word line drive circuit.\nAn intermediate potential generation circuit will first be described. An example of an intermediate potential generation circuit is disclosed in IEEE Journal of Solid-State Circuit, Vol. SC-22, No. 5, October 1987, pp. 861-867. FIG. 25 is a circuit diagram showing a structure of a conventional intermediate potential generation circuit thereof.\nReferring to FIG. 25, an intermediate potential generation circuit includes transistors Q101-Q103 which are n type MOSFETS, transistors Q104-Q106 which are p type MOSFETs, and resistors R101-R104.\nFIG. 26 schematically shows a structure of the intermediate potential generation circuit of FIG. 25 on a p type substrate. Referring to FIG. 26, the intermediate potential generation circuit includes a p type substrate 111, an n type well 112, transistors Q11-Q106, and resistors R101-R104. In FIG. 26, components corresponding to those of FIG. 25 have the same reference characters denoted.\nAn operation of the intermediate potential generation circuit will be described hereinafter with reference to FIGS. 25 and 26.\nThe resistance of resistors R101 and R102 equal each other. Also, the resistance of resistors R103 and R104 equal each other. The resistance of resistors R101-R104 is several m.OMEGA., which is high resistance. Therefore, the current flowing in transistors Q101, Q102, Q104 and Q105 is reduced, and these transistors conduct lightly. Therefore, the gate-source potential of transistor Q101, Q102, Q104 and Q105 is equal to the threshold voltage of each transistor.\nAccording to the above-described structure, the potential of nodes N1 and N3 is approximately V.sub.CC /2 (V.sub.CC is the power supply voltage). Therefore, the potential of node N2 becomes V.sub.CC /2+V.sub.TH101 (V.sub.TH101 is the threshold voltage of transistor Q101), and the potential of node N4 is approximately V.sub.CC /2-.vertline.V.sub.TH105 .vertline. (V.sub.TH105 is the threshold voltage of transistor Q105). When the potential of an output signal V.sub.sg is lower than V.sub.CC /2+V.sub.TH101 -V.sub.TH103 (V.sub.TH103 is the threshold voltage of transistor Q103), transistor Q103 conducts, whereby the potential of output signal V.sub.sg rises. When the potential of output signal V.sub.sg is higher than V.sub.CC /2-.vertline.V.sub.TH105 .vertline.+.vertline.V.sub.TH106 .vertline. (V.sub.TH106 is the threshold voltage of transistor Q106), transistor Q106 conducts, whereby the potential of output signal V.sub.sg falls. By the above-described operation, the potential of output signal V.sub.sg becomes approximately V.sub.CC /2.\nA timer circuit for a self refresh operation will be described hereinafter. A refresh operation must be carried out periodically since a DRAM is a volatile memory. Lengthening the period of a refresh operation will reduce power consumption thereof, to allow reduction of power consumption in the device. In a conventional timer circuit, a refresh operation is carried out when the potential held in a memory cell becomes lower than a predetermined level. An example of such a timer circuit is disclosed in IEEE Journal of Solid-State Circuits, Vol. 26, No. 11, November 1991, pp. 1556-1562. FIG. 27 shows a structure of this conventional timer circuit.\nReferring to FIG. 27, a timer circuit includes a differential amplifier 121, an S-R flipflop 122, a delay circuit 123, a transistor Q111 which is an n type MOSFET, a capacitor 124 of a memory cell, and an n type diffusion layer 125.\nFIGS. 28A(1) to 28B(B) are timing chart showing the operation of the timer circuit of FIG. 27.\nAn operation of the timer circuit will be described hereinafter with reference to FIGS. 27 and 28A(1) to 28B(B). When the potential V.sub.N in capacitor 124 becomes lower than a reference potential V.sub.REF at time t.sub.1, S-R flipflop 122 is set to render the level of an output signal .phi..sub.E to a H level (logical high). Output signal .phi..sub.E of S-R flipflop 122 is delayed for a predetermined time, and then applied to a reset terminal R of S-R flipflop 122. As a result, a reset signal R attains a H level. This causes S-R flipflop 122 to be reset, whereby output signal .phi..sub.E attains a L level (logical low). A refresh operation is carried while output signal .phi..sub.E attains a H level, whereby transistor Q111 attains a conductive state, and the potential of capacitor 124 of a memory cell is maintained at V.sub.CC. Then, when output signal .phi..sub.E attains a L level, transistor Q111 is rendered non-conductive, whereby the holding voltage V.sub.N of capacitor 124 is gradually reduced by leakage current. When holding voltage V.sub.N of capacitor 124 becomes lower than reference voltage V.sub.REF, an operation similar to that of the above-described operation is repeated. Thus, a refresh operation is carried out at a predetermined period.\nAn internal high voltage circuit will be described. FIG. 29 is a block diagram showing a structure of a conventional internal high voltage circuit. Referring to FIG. 29, an internal high voltage circuit includes a first detector 132, a second detector 132, a third detector 133, a first oscillator 134, a second oscillator 135, a small pump 136, a large pump 137, a RAS pump 138, and an AND gate G101 and an inverter G102.\nWhen high voltage V.sub.PP supplied to a word line driver 139 becomes lower than a predetermined potential, first detector 131 provides an output signal .phi..sub.E1 of a H level to first oscillator 134. First oscillator 134 oscillates while output signal .phi..sub.E1, attains a H level, and provides an oscillation signal to small pump 136. Small pump 136 responds to this oscillation signal to provide high voltage V.sub.PP to word line driver 139 at a standby state.\nWhen the high voltage supplied to word line driver 139 becomes lower than a predetermined potential, second detector 132 provides an output signal .phi..sub.E2 of H level to second oscillator 135. Second oscillator 135 oscillates when output signal .phi..sub.E2 attains a H level, and provides an oscillation signal to large pump 137. Large pump 137 responds to this oscillation signal to rapidly increase high voltage V.sub.PP supplied to word line driver 139.\nWhen high voltage V.sub.PP supplied to word line driver 139 becomes lower than a predetermined potential, third detector 133 provides an output signal .phi..sub.E3 of a H level to AND gate G101. AND gate G101 takes the logical product of output signal .phi..sub.E3 and an inverted signal of a row address strobe signal /RAS (\"/\" implies a low-active signal) to provide an output signal to RAS pump 138. AND gate G101 provides an output signal when row address strobe signal /RAS attains a L level, whereby the semiconductor device operates to raise the word line to high voltage V.sub.PP.\nThe first detector shown in FIG. 29 will be described with reference to FIG. 30 showing a circuit diagram thereof.\nReferring to FIG. 30, a first detector includes transistors Q121-Q124 which are p type MOSFETs, and transistors Q125 and Q126 which are n type MOSFETs.\nHigh voltage V.sub.PP provided to the first detector is reduced by a threshold voltage .sub.TH of each transistor, i.e., by 3V.sup.TH, by transistors Q121-Q123. Therefore, an output signal .phi..sub.E1 of a H level is output when high voltage V.sub.PP becomes lower than V.sub.CC +3V.sub.TH.\nThe second detector of FIG. 29 will be described hereinafter with reference to FIG. 31 showing a circuit diagram thereof.\nReferring to FIG. 31, a second detector includes transistors Q131-Q133 which are p type MOSFETs, and transistors Q134 and Q135 which are n type MOSFETs.\nHigh voltage V.sub.PP provided to the second detector is reduced by a threshold voltage V.sub.TH of each transistor, i.e. 2V.sup.TH, by transistors Q131 and Q132. Therefore, the second detector provides an output signal .phi..sub.E2 of a H level when high voltage V.sub.PP becomes lower than V.sub.CC +2V.sub.TH. The third detector of FIG. 29 has a structure similar to that of the second detector of FIG. 31, and also operates in a similar manner thereof.\nThe first oscillator of FIG. 29 will be described hereinafter with reference to FIG. 32 showing a circuit diagram thereof.\nReferring to FIG. 32, a first oscillator includes transistors Q141-Q148 which are p type MOSFETs, and transistors Q149-Q156 which are n type MOSFETs. C101-C103 shown in FIG. 20 are the parasitic capacitance of each portion.\nBecause transistor Q141 has a long channel length, current flowing in transistor Q149 is limited to a current value of I.sub.1. Transistor Q149 and transistors Q150, Q152, Q154 and Q156 form a current mirror, so that current flowing through transistors Q143, Q145, Q147, Q152, Q154 and Q156 is limited to the value of I.sub.1. Therefore, the delay time of each inverter formed by each of these transistors becomes 3C/I.sub.1 where each capacitance of parasitic capacitances C101-C103 is C.\nWhen V.sub.CC /2+V.sub.TH101 -V.sub.TH103 >V.sub.sg =V.sub.CC /2, and V.sub.CC /2-.vertline.V.sub.TH105 .vertline.+.vertline.V.sub.TH106 .vertline.<V.sub.sg =V.sub.CC /2, i.e. V.sub.TH101 >V.sub.H103, and .vertline.V.sub.TH105 .vertline.>.vertline.V.sub.TH106 .vertline. in the intermediate potential generation circuit of FIG. 25, a through current flows in transistors Q103 and Q106 at the time of standby since transistors Q103 and Q106 both conduct when the potential of output signal V.sub.sg is stable at V.sub.CC /2. There was a problem that the power consumption of the device was increased due to this through current.\nIn the timer circuit of FIG. 27, the period of a refresh operation is T.sub.1, at low temperature as shown in FIG. 28(a), and is T.sub.2 at high temperature as shown in FIG. 28(b) because leakage current of capacitor 124 increases at high temperature.\nThe timer circuit shown in FIG. 27 had problems set forth in the following. A phenomenon called soft error is seen in a DRAM. More specifically, .alpha. particles emitted from the package or the like cause the generated electrons to be captured in an n type diffusion layer 125 of a memory cell, whereby information in the memory cell is destroyed. Therefore, soft error easily occurs when the holding voltage V.sub.N becomes not higher than the lowest holding voltage V.sub.REF required for proper operation of a readout circuit of a memory cell by a predetermined value of .DELTA.V. As a result, when the level of holding voltage V.sub.REF is equal at both the high and low temperature, the time period having a high probability of generating soft error is d.sub.1, and d.sub.2 at a low temperature and a high temperature, respectively, as shown in FIGS. 28A and 28B. Therefore, there was a problem that the possibility of soft error occurrence is increased at low temperature.\nIn the first and second detectors shown in FIGS. 30 and 31, a through current is conducted to increase power consumption since all transistors Q124, Q126, Q133 and Q135 become conductive when the level of output signals .phi..sup.E1 and .phi..sub.E2 change.\nIn the third detector shown in FIG. 31, the time required for pulling the potential of the node between transistorS Q132 and Q134 from a H level to a L level is several .mu.s. An operation of a DRAM occurs at the minimum of every 90 ns, for example. Therefore, a word line is driven several ten times during the transition of the third director from an off state to an on state, resulting in reduction of the level of high voltage V.sub.PP of word line driver 139. FIGS. 33(a) to 33(b) are diagrams for describing the change in the level of high voltage V.sub.PP, with respect to output signal .phi..sub.E3 of the third director. It is appreciated from FIG. 33(a) to 33(b) that the level of high voltage V.sub.PP is gradually reduced according to each transition of row address strobe signal /RAS when output signal .phi..sub.E3 attains a L level. Therefore, a conventional third detector is set so that sufficient current is conducted to transistor Q134 in order to rapidly pull down the potential of the node between transistors Q132 and Q134 rapidly to a L level from a H level. Thus, there was a problem that power consumption is increased during standby.\nIn the first oscillator of FIG. 32, the delay time of 3C/I.sub.1 is reduced due to increase of current I.sub.1 flowing in transistor Q141 in response to increase of power supply potential V.sub.CC. This causes the oscillation frequency of the first oscillator to be increased to shorten the operation cycle. Thus, there was a problem that power consumption of the device is increased."}
-{"text": "This invention relates generally to coin handling and processing apparatus such as coin wrapping apparatus and more particularly to a coin stacking tube device suitable for incorporation in any of such apparatus.\nA typical example of a conventional coin wrapping apparatus is that of a construction wherein coins supplied from a hopper onto a turntable are arrayed by centrifugal force applied thereto along the periphery of the turntable, and a coin counting mechanism counts the coins while the coins are sent from the turntable to a coin passage and propelled therealong by, for instance, a propelling belt. The coins delivered from the coin passage one after another are received in a coin stacking tube to be stacked therein, the coins thus stacked being dropped into a coin wrapping device comprising a plurality of coin wrapping rolls disposed around a circle directly below the coin stacking tube. The thus stacked coins are wrapped by the wrapping rolls with a piece of paper, and the lateral edges of the paper are fold crimped to form firm beads by which the paper is maintained in tightly wrapped state.\nIn the case when it is required to change the denomination of coins to be wrapped in the coin wrapping apparatus, various parts of the apparatus must be readjusted so that the apparatus is set for the new denomination of coins. Of these parts, the number of coins supplied onto the turntable, the rotating speed of the same, the lateral width of the coin passage, the height and the position of the propelling belt, the positions of the wrapping rolls to be brought into contact with the stack of coins, and the rotating speed of the wrapping rolls can be readjusted comparatively easily by interlinking those members controlling these values with a member for setting the apparatus to a different denominations of coins.\nHowever, the adjustment of the inner diameter of the coin stacking tube is not easy, and therefore it has been a conventional practice to prepare a number of coin stacking tubes each having a different inner diameter suitable for a specific denomination (or one outer diameter) of coins and, at the time of changing the denomination of coins to be wrapped, to replace the existing coin stacking tube with another suitable for the new denomination.\nIn this case also, the denomination changing operation is found to be troublesome because of the requirement of the selection and the replacement, and furthermore there has been a high possibility of erroneous selection of the coin stacking tube, causing unsatisfactory stacking of coins in the erroneous coin stacking tube.\nRecently, a type of coin stacking tube whose inner diameter is made variable in accordance with the denomination of coins has been developed. In this kind of coin stacking tube, the space for receiving the coins is formed by a plurality of blades provided in an outer casing, the inner edge of each blace contacting against the inner surface of the preceding blade, the stem part of each blade being supported rotatably about an axis fixed to the outer casing, and all blades being rotated in either of the opening and closing directions in accordance with the denomination setting in the coin wrapping apparatus, whereby the space formed within the plurality of blades is adapted for the denomination or the outer diameter of the coins to be stacked therein.\nIn this example, however, in order to assure smooth rotation of the blades at the time of expension or constraction of the interior space of the coin stacking tube, the rotating axes for the blades must be secured to the outer casing in parallel with each other. If the rotating axes are not accurately, parallel, the coin stacking space formed therein cannot be of a correct circular cross-section, thus causing unsatisfactory stacking of coins within the tube, and the uneven contacting of the innermost edges of the blades against the inner surfaces of the preceding blades causing a considerable torque to be required for rotating the blades."}
-{"text": "The present invention is directed to a fan array fan section utilized in an air-handling system.\nAir-handling systems (also referred to as an air handler) have traditionally been used to condition buildings or rooms (hereinafter referred to as \u201cstructures\u201d). An air-handling system is defined as a system that includes components designed to work together in order to condition air as part of the primary system for ventilation of structures. The air-handling system may contain components such as cooling coils, heating coils, filters, humidifiers, fans, sound attenuators, controls, and other devices functioning to meet the needs of the structures. The air-handling system may be manufactured in a factory and brought to the structure to be installed or it may be built on site using the necessary devices to meet the functioning needs of the structure. The air-handling compartment 102 of the air-handling system includes the inlet plenum 112 prior to the fan inlet cone 104 and the discharge plenum 110. Within the air-handling compartment 102 is situated the fan unit 100 (shown in FIGS. 1 and 2 as an inlet cone 104, a fan 106, and a motor 108), fan frame, and any appurtenance associated with the function of the fan (e.g. dampers, controls, settling means, and associated cabinetry). Within the fan 106 is a fan wheel (not shown) having at least one blade. The fan wheel has a fan wheel diameter that is measured from one side of the outer periphery of the fan wheel to the opposite side of the outer periphery of the fan wheel. The dimensions of the handling compartment 102 such as height, width, and airway length are determined by consulting fan manufacturers data for the type of fan selected.\nFIG. 1 shows an exemplary prior art air-handling system having a single fan unit 100 housed in an air-handling compartment 102. For exemplary purposes, the fan unit 100 is shown having an inlet cone 104, a fan 106, and a motor 108. Larger structures, structures requiring greater air volume, or structures requiring higher or lower temperatures have generally needed a larger fan unit 100 and a generally correspondingly larger air-handling compartment 102.\nAs shown in FIG. 1, an air-handling compartment 102 is substantially divided into a discharge plenum 110 and an inlet plenum 112. The combined discharge plenum 110 and the inlet plenum 112 can be referred to as the airway path 120. The fan unit 100 may be situated in the discharge plenum 110 as shown), the inlet plenum 112, or partially within the inlet plenum 112 and partially within the discharge plenum 110. The portion of the airway path 120 in which the fan unit 100 is positioned may be generically referred to as the \u201cfan section\u201d (indicated by reference numeral 114). The size of the inlet cone 104, the size of the fan 106, the size the motor 108, and the size of the fan frame (not shown) at least partially determine the length of the airway path 120. Filter banks 122 and/or cooling coils (not shown) may be added to the system either upstream or downstream of the fan units 100.\nFor example, a first exemplary structure requiring 50,000 cubic feet per minute of air flow at six (6) inches water gage pressure would generally require a prior art air-handling compartment 102 large enough to house a 55 inch impeller, a 100 horsepower motor, and supporting framework. The prior art air-handling compartment 102, in turn would be approximately 92 inches high by 114 to 147 inches wide and 106 to 112 inches long. The minimum length of the air-handling compartment 102 and/or airway path 120 would be dictated by published manufacturers data for a given fan type, motor size, and application. Prior art cabinet sizing guides show exemplary rules for configuring an air-handling compartment 102. These rules are based on optimization, regulations, and experimentation.\nFor example, a second exemplary structure includes a recirculation air handler used in semiconductor and pharmaceutical clean rooms requiring 26,000 cubic feet per minute at two (2) inches water gage pressure. This structure would generally require a prior art air-handling system with a air-handling compartment 102 large enough to house a 44 inch impeller, a 25 horsepower motor, and supporting framework. The prior art air-handling compartment 102, in turn would be approximately 78 inches high by 99 inches wide and 94 to 100 inches long. The minimum length of the air-handling compartment 102 and/or airway path 120 would be dictated by published manufacturers data for a given fan type, motor size and application. Prior art cabinet sizing guides show exemplary rules for configuring an air-handling compartment 102. These rules are based on optimization, regulations, and experimentation.\nThese prior art air-handling systems have many problems including the following exemplary problems:                Because real estate (e.g. structure space) is extremely expensive, the larger size of the air-handling compartment 102 is extremely undesirable.        The single fan units 100 are expensive to produce and are generally custom produced for each job.        Single fan units 100 are expensive to operate.        Single fan units 100 are inefficient in that they only have optimal or peak efficiency over a small portion of their operating range.        If a single fan unit 100 breaks down, there is no air conditioning at all.        The low frequency sound of the large fan unit 100 is hard to attenuate.        The high mass and turbulence of the large fan unit 100 can cause undesirable vibration.        \nHeight restrictions have necessitated the use of air-handling systems built with two fan units 100 arranged horizontally adjacent to each other. It should be noted, however, that a good engineering practice is to design air handler cabinets and discharge plenums 110 to be symmetrical to facilitate more uniform air flow across the width and height of the cabinet. Twin fan units 100 have been utilized where there is a height restriction and the unit is designed with a high aspect ratio to accommodate the desired flow rate. As shown in the Greenheck \u201cInstallation Operating and Maintenance Manual,\u201d if side-by-side installation was contemplated, there were specific instructions to arrange the fans such that there was at least one fan wheel diameter spacing between the fan wheels and at least one-half a fan wheel diameter between the fan and the walls or ceilings. The Greenheck reference even specifically states that arrangements \u201cwith less spacing will experience performance losses.\u201d Normally, the air-handling system and air-handling compartment 102 are designed for a uniform velocity gradient of 500 feet per minute velocity in the direction of air flow. The two fan unit 100 air-handling systems, however, still substantially suffered from the problems of the single unit embodiments. There was no recognition of advantages by increasing the number of fan units 100 from one to two. Further, the two fan unit 100 section exhibits a non-uniform velocity gradient in the region following the fan unit 100 that creates uneven air flow across filters, coils, and sound attenuators.\nIt should be noted that electrical devices have taken advantage of multiple fan cooling systems. For example, U.S. Pat. No. 6,414,845 to Bonet uses a multiple-fan modular cooling component for installation in multiple component-bay electronic devices. Although some of the advantages realized in the Bonet system would be realized in the present system, there are significant differences. For example, the Bonet system is designed to facilitate electronic component cooling by directing the output from each fan to a specific device or area. The Bonet system would not work to direct air flow to all devices in the direction of general air flow. Other patents such as U.S. Pat. No. 4,767,262 to Simon and U.S. Pat. No. 6,388,880 to El-Ghobashy et al. teach fan arrays for use with electronics.\nEven in the computer and machine industries, however, operating fans in parallel is taught against as not providing the desired results except in low system resistance situations where fans operate in near free delivery. For example, Sunon Group has a web page in which they show two axial fans operating in parallel, but specifically state that if \u201cthe parallel fans are applied to the higher system resistance that [an] enclosure has, . . . less increase in flow results with parallel fan operation.\u201d Similar examples of teaching against using fans in parallel are found in an article accessible from HighBeam Research's library (http://stati.highbeam.com) and an article by Ian McLeod accessible at (http://www.papstplc.com)."}
-{"text": "Adjustable water ski bindings are utilized to attach water skis to the feet of skiers and consist of a foot piece which is stationary on the water ski and a heel piece which is adjustable lengthwise of the water ski. In order to attach the water ski to the foot, the toe is first inserted into the front vamp or foot piece and then the heel piece is pressed against the heel until a tight fit on the foot is accomplished. The heel piece can be attached to a mounting plate which is slidable in guides attached to the ski on opposite sides of the plate. The mounting plate can carry two latch members having teeth which cooperate with ratchet teeth on the guides and when the teeth are engaged, the heel piece will be locked against rearward movement which would loosen the heel piece. Quick loosening of the heel piece can be accomplished by disengagement of the teeth by movement of the latch members so that the skier is able to quickly remove the ski from his foot in the event of a fall into the water. Present locking mechanisms for the heel piece utilize pawl or latch members and springs which are unnecessarily complicated since either pivots or guide slots are required in the mounting plate for the latch members and separate springs are required for these members. The U.S. Pat. to W. J. Meucci No. 3,137,014 is an example of guide slots cut in the mounting plate for latch or pawl members. The use of pivots connected to the mounting plate for latch or pawl members is illustrated by U.S. Pat. Nos. to H. A. Moline 2,970,325; B Roudebush 3,127,623; and W. W. Bennett 3,102,279. Also, the patent to R. I. Rumig, 2,866,210 requires a separate spring attached to the mounting plate for the latch member. These prior locking mechanisms for ski bindings are unnecessarily expensive and are, under some circumstances, difficult to operate during attachment and removal of the ski from the foot."}
-{"text": "Conventionally, as a means of displaying images, projection-type optical display apparatuses such as projectors are known. Such optical display apparatuses require an optical illumination apparatus for efficiently and uniformly illuminating the optical image formed on a display panel, such as a reflective liquid crystal display panel. FIG. 10 is a diagram conceptually showing an example of the construction of an optical display apparatus employing a conventional optical illumination apparatus.\nIn FIG. 10, reference numeral 101 represents a light source, and reference numeral 102 represents a reflector disposed so as to partially surround the light source 101. A PBS (polarizing beam splitter) prism unit 103 is disposed immediately behind the reflector 102, i.e., on the right side thereof in FIG. 10. The PBS prism unit 103 includes a plurality of PBS prisms arranged parallel to one another. The PBS prism unit 103 splits the light from the light source 101 into two differently polarized types of light. Of the individual PBS prisms 103a and 103b, those which let out S-polarized light as described later have half-wave plates 104 disposed immediately behind them.\nBehind the PBS prism unit 103 (i.e., on the right side thereof in FIG. 10) are disposed, in order of arrangement, a first lens array 105, then somewhat away therefrom, a second lens array 106, and a superimposing lens 107 immediately behind it. The first lens array 105 has a plurality of lens cells 105a arranged in a rectangular, grid-like array having an aspect ratio substantially identical to that of a display panel 109 to be described later. Similarly, the second lens array 106 also has a plurality of lens cells 106a arranged in a rectangular, grid-like array. However, the shape of the lens cells 106a of the second lens array 106 is not necessarily geometrically similar to that of the lens cells 105a. \nThe images from the individual lens cells 105a of the first lens array 105 are, by the second lens array 106 and the superimposing lens 107 disposed immediately behind it, superimposed on one another in the vicinity of the focal point of the superimposing lens 107. The display panel 109 is disposed at the focal point of the superimposing lens 107. The display panel 109 is illuminated in a telecentric fashion by a condenser lens 108 disposed immediately in front of it. The components from the first lens array 105 through the superimposing lens 107 mentioned above together constitute an optical integrator system. It is to be noted that, in all the diagrams referred to in the present specification, irrespective of whether they relate to prior-art examples or to embodiments of the present invention, light beams are represented by their optical axes alone.\nThe light emitted from the light source 101 is reflected from the reflector 102, and is thereby formed into a substantially parallel beam and directed to the PBS prisms 103a of the PBS prism unit 103. Here, P-polarized light, indicated by solid lines P, is transmitted straight through the PBS prisms 103a. On the other hand, S-polarized light, indicated by broken lines S, is reflected inside the PBS prisms 103a so as to be directed to the outwardly contiguous PBS prisms 103b, and is then reflected again inside the PBS prisms 103b so as to exit therefrom, still as S-polarized light. That is, by the PBS prism unit 103, the light from the light source 101 is split into two differently polarized types of light in the direction of the longer sides of the display panel 109, i.e., in a vertical direction along the plane of the figure.\nThe S-polarized light exiting from the PBS prisms 103b is transmitted through the half-wave plates 104 disposed immediately behind the PBS prisms 103b and is thereby converted into P-polarized light. That is, a portion of the light from the light source 101 has its polarization converted first by the PBS prisms 103b of the PBS prism unit 103 and then by the half-wave plates 104, and eventually comes out as uniformly P-polarized light. This arrangement constitutes a polarization conversion device. Here, the type of light into which the light from the light source 101 is converted does not necessarily have to be P-polarized light, but can be of other polarizations. The arrangement described thus far, starting with the light source 101 and ending immediately in front of the display panel 109, constitutes an optical illumination apparatus.\nThe light thus converted into uniformly P-polarized light is then directed through the above-mentioned optical integrator system to the display panel 109. The display panel 109 modulates, pixel by pixel, the light it is illuminated with according to the display data fed thereto, and emits the modulated light. The light thus emitted then enters an optical projection system 110. The display data presented on the display panel 109 is projected, as an image, onto a screen (not shown) through this optical projection system 110. Reference numeral 110a represents an aperture stop disposed in the optical projection system 110.\nFIG. 11 is a diagram conceptually showing another example of the construction of an optical display apparatus employing a conventional optical illumination apparatus. In this figure, reference numeral 201 represents a light source, and reference numeral 202 represents a reflector disposed so as to partially surround the light source 201. Behind the reflector 202 (i.e., on the right side thereof in FIG. 11) are disposed, in order of arrangement, a first lens array 203 and, then somewhat away therefrom, a second lens array 204. The first lens array 203 has a plurality of lens cells 203a arranged in a rectangular, grid-like array having an aspect ratio substantially identical to that of a display panel 209 to be described later. Similarly, the second lens array 204 also has a plurality of lens cells 204a arranged in a rectangular, grid-like array. However, the shape of the lens cells 204a of the second lens array 204 is not necessarily geometrically similar to that of the lens cells 203a. \nA PBS (polarizing beam splitter) prism array 205 is disposed immediately behind the second lens array 204. The PBS prism array 205 includes a plurality of PBS prisms arranged in an array. The PBS prism array 205 splits the light from the light source 201 into two differently polarized types of light. Of the individual PBS prisms 205a and 205b, those which let out S-polarized light as described later have half-wave plates 206 disposed immediately behind them.\nA superimposing lens 207 is disposed behind the PBS prism array 205. The images of the individual lens cells 203a of the first lens array 203 are, by the second lens array 204 and the superimposing lens 207, superimposed on one another in the vicinity of the focal point of the superimposing lens 207. The display panel 209 is disposed at the focal point of the superimposing lens 207. The display panel 209 is illuminated in a telecentric fashion by a condenser lens 208 disposed immediately in front of it. The first lens array 203, the second lens array 204, and the superimposing lens 207 mentioned above together constitute an optical integrator system.\nThe light emitted from the light source 201 is reflected from the reflector 202, and is thereby formed into a substantially parallel beam and passed through the first lens array 203 and the second lens array 204, so that the light exiting from the individual lens cells 204a of the second lens array 204 enters corresponding ones of the PBS prisms 205a of the PBS prism array 205. Here, P-polarized light, indicated by solid lines P, is transmitted straight through the PBS prisms 205a. On the other hand, S-polarized light, indicated by broken lines S, is reflected inside the PBS prisms 205a so as to be directed to the contiguous PBS prisms 205b, and is then reflected again inside the PBS prisms 205b so as to exit therefrom, still as S-polarized light.\nThe S-polarized light exiting from the PBS prisms 205b is then transmitted through the half-wave plates 206 disposed immediately behind the PBS prisms 205b and is thereby converted into P-polarized light. That is, a portion of the light from the light source 201 has its polarization converted first by the PBS prisms 205b of the PBS prism array 205 and then by the half-wave plates 206, and eventually comes out as uniformly P-polarized light. This arrangement constitutes a polarization conversion device. Here, the type of light into which the light from the light source 201 is converted does not necessarily have to be P-polarized light, but can be of other polarizations. The arrangement described thus far, starting with the light source 201 and ending immediately in front of the display panel 209, constitutes an optical illumination apparatus.\nThe light thus converted into uniformly P-polarized light is then directed through the superimposing lens 207 to the display panel 209. The display panel 209 modulates, pixel by pixel, the light it is illuminated with according to the display data fed thereto, and emits the modulated light. The light thus emitted then enters an optical projection system 210. The display data presented on the display panel 209 is projected, as an image, onto a screen (not shown) through this optical projection system 210. Reference numeral 210a represents an aperture stop disposed in the optical projection system 210.\nIn the conventional optical illumination apparatus constructed as shown in FIG. 10, polarization conversion is performed immediately behind the light source 101. Therefore, the light emitted from the light source 101 and then reflected from the reflector 102 has its beam diameter enlarged to about twice its original beam diameter as a result of the polarization conversion. This diminishes the f-number of the illumination light Ia that strikes the display panel 109 and thus diminishes the f-number of the projection light Ea that emanates from the display panel 109, making the burden on the optical projection system 110 heavier.\nOn the other hand, in the conventional optical illumination apparatus constructed as shown in FIG. 11, the light emitted from the light source 201 and then reflected from the reflector 202 experiences no enlargement of its beam diameter. Therefore, no diminishing occurs in the f-number of the illumination light Ib that strikes the display panel 209 nor in the f-number of the projection light Eb that emanates from the display panel 209. Thus, no extra burden is placed on the optical projection system 210. However, the light from the light source 201 is not converted into uniformly polarized light until it has passed through the second lens array 204. Therefore, in this optical illumination apparatus, no space is available for inserting a polarization-dependent color switching device such as those used in the embodiments of the present invention to be described later. That is, this optical illumination apparatus does not permit a so-called color sequential illumination method using such a color switching device."}
-{"text": "Certain presentations and graphs require printing on a printable media, such as a paper sheet, of a dimension that is most beneficially presented in the form of a strip, for example, presentation charts used in project management. Project management is the planning and control of many activities that must be coordinated to achieve specific goals leading to the completion of an overall given objective. The project management process frequently uses a set of tools which incorporate charts and reports to detail the project for communication within the project team and with others.\nOne form of project management chart is referred to as a work breakdown structure (WBS) chart. A WBS chart is an organizational diagram type of chart depicting work packages comprising all of the principal elements of a project. Another form of project management chart, used for communication with the project team and with others, is a precedent network (Network) chart, sometimes erroneously referred to as a PERT chart. The Network chart displays activities required to produce the work packages depicted in the WBS chart and shows the relationships between the activities, i.e. the precedents and dependencies between the activities as they flow towards completion of a project.\nFor practical project management purposes, the WBS and Network charts are usually more than one standard sized page in width. To present either chart may require many standard sized pages to be joined one to the other. As the project evolves, the project management process results in the updating, modifying, and reproducing the project charts as a consequence of project progress. Thus, the WBS and network charts will change during the course of the project, consequently requiring the WBS and network charts to be produced frequently during the course of the project. To allow the charts to be printed on standard sized paper using conventional computer printers or photocopied onto standard sized paper, requires the joining of standard sized pages together to form a completed WBS or Network chart. Joining these pages is a time consuming process and usually requires input from the project manager to lay out the pages in a proper sequence prior to cutting and pasting them together to form the chart. Each time revised charts are produced during the life of a project, several copies are required each for key team members. The page assembly process to produce the charts is a frustrating task and often results in sloppy presentations even though the computer-generated data or the images printed or photocopied on each of the individual pages may be perfect."}
-{"text": "The primary source for oxygen for a grill comes from below the burner. Conventional grills have vents located near the bottom of the grill, below the burner, to allow for air flow which supplies the necessary oxygen. However, if the airflow through these vents is at too high a velocity, the flame at the burner can be blown out. This can cause conventional grills to be virtually useless. in high wind environments.\nAdditionally, other areas of high velocity airflow can also disrupt the burner operation. One such area is the venturi assembly. Conventional grills have a venturi tube for mixing propane and air which is generally located outside of the grill tub. High wind conditions can result in this mixture being diluted too much, resulting in the burner flame going out. Additionally, high velocity airflow in the area directly above the burner can disrupt the burner operation. This can occur either when the grill hood is open, allowing the wind to access the top of the burner through the heat diffusing material, which typically consists of rocks or ceramic bricketts. Wind can also affect burner operation through vents in the grill hood when the hood is closed. In any of these cases, the high wind conditions can cause the flame at the burner to go out, and thus the grill must be re-lit, if possible, in order to continue operation.\nBased on the foregoing, there is a need for a grill which overcomes these problems and protects the flame at the burner from high velocity airflow."}
-{"text": "Hot beverage makers, such as coffee makers, have been known and sold for many years using various brewing techniques. The typical and traditional coffee maker includes a stand or tower that has a warming plate forming the bottom or base of the tower with a filter basket located above the warming plate. The interior of the tower defines, at least in part, a fresh water reservoir. Such coffee makers further include a fluid reservoir, such as a glass carafe, that rests on the warming plate beneath the filter basket. Alternatively, the fluid reservoir may be an insulated carafe, in which case the warming plate is typically omitted.\nIn use, an operator fills the carafe in order to transfer water to the fresh water reservoir. The water is heated and passed through the filter basket, which includes the grounds to be infused. The brewed beverage then flows from the basket into the carafe. The beverage is maintained at an elevated temperature via the warming plate upon which the carafe rests, in the case of a glass carafe, or by the insulating properties of the carafe in the case of an insulated carafe.\nA new variation of coffee maker has been developed wherein a brewed beverage tank is included, such as the coffee makers described in commonly assigned U.S. Pat. No. 6,564,975 to Garman, issued May 20, 2003, and U.S. Pat. No. 6,681,960 to Garman, issued Jan. 27, 2004, the contents of which are incorporated herein in their entirety. Briefly, the brewed beverage tank holds the filter basket above a reservoir portion. Hot water passes through the filter basket and a material to be infused (e.g., ground coffee beans). The brewed beverage is then collected and held in the reservoir portion of the brewed beverage tank (\u201cbrew tank\u201d). A dispenser actuator is depressed that opens an outlet port in the reservoir. A user simply actuates the dispenser actuator with a mug or cup and the brewed beverage passes through the outlet to the operator's container.\nIn order to allow a user to determine the amount of brewed coffee within the device, a transparent window on the coffee holding reservoir can be used. A transparent vertical slot-type window on the side of the holding reservoir can also be used. Also known is the use of a buoyant ball contained within a transparent column, the transparent column containing coffee at the same level as in the holding reservoir and the ball floating on the surface of the coffee in the column. The transparent column typically includes markings that correspond to the quantity of coffee in the reservoir, such that the quantity of coffee in the reservoir may be determined by reading the marking nearest the floating ball.\nHowever, these transparent windows, slots, and columns give a somewhat rough estimate of how much coffee remains and often do not give a clear indication of the exact amount. In addition, the amount of coffee remaining in the coffee maker is often difficult to gauge as the coffee may be difficult to see against the background of the coffee maker which is often somewhat similar in color to the coffee itself. Further, the glass or plastic of the these transparent windows, slots, and columns may be quickly stained by the coffee, such that the transparency is greatly reduced and the amount of coffee becomes difficult or impossible to determine. Because of the above problems, even when the amount of coffee may be determined at a short distance from the coffee maker via these transparent windows, slots, and columns, it may be difficult or impossible to determine the amount of coffee at a greater distance, such as across a room.\nThere is a need, therefore, for a coffee level indicator that is accurate, reliable, clearly indicates the amount of brewed coffee remaining in the holding reservoir, is unaffected by staining, and is readable at a distance. The operation and structure of a brew tank with an integrated fluid gauge in accordance with the present invention would solve one or more of these or other needs."}
-{"text": "1. Field of the Invention\nThe present invention relates to an image recording apparatus in which is mounted a recording head that performs recording by ejecting (discharging) a liquid from an energy generating element or by thermal transfer.\nThe present invention can be applied for apparatuses, such as printers, copiers, facsimile machines for which communication systems are provided, or word processors that incorporate printers, that perform the recording of images on a recording medium, such as paper, thread, fiber, cloth, leather, metal, plastic, glass, wood or ceramics, and for industrial recording apparatuses with which various processors are combined.\nxe2x80x9cRecordingxe2x80x9d in this invention is defined not only as the formation on a recording medium of images, such as characters or drawings, that convey meaning, but also as the formation of images, such as patterns, that convey no meaning.\n2. Related Background Art\nConventionally, the demand for recording apparatuses that can produce high quality images has increased, and how to improve image quality has been the subject of numerous discussions. For a recording apparatus in which a recording head is moved in one direction when recording images, the precision of the positioning of an image to be recorded is determined by the accuracy with which the recording head itself is positioned. And for the improvement of the image quality, the enhancement of the accuracy with which a recording head is positioned is an extremely important element. Therefore, in a conventional recording apparatus, for a carriage on which is mounted a recording head that records in only one direction, position detection means (e.g., an image scanner) is provided for accurately ascertaining the position of the recording head. Or, at the carriage\"\"s home position in the apparatus, optical reading means is provided to detect the position of the recording head. Then, based-on the obtained head positioning data, whether the recording position is adequate or whether the recording position must be corrected is determined.\nHowever, in a conventional recording apparatus the recording head, which constitutes the printing means, and the position detection means are arranged separately. Therefore, in a recording apparatus wherein, for example, a head position detection means is provided for a carriage, satisfactory positioning accuracy for the recording head must be obtained by mounting the recording head on the carriage. In order to obtain such accuracy, precision in the sizing of components, such as the carriage and the recording head, must be improved, or a process must be performed for correcting the positioning of the recording head.\nIn addition, since elements and circuits for detecting the position of the recording head must be formed on the carriage or on the substrate of the apparatus, manufacturing costs will be increased.\nFrom the viewpoint of high quality image recording, highly delicate recording, for improved image density and tone representation, can be performed by producing dots that have variable sizes.\nAs the resolution of an image is increased, however, extremely high accuracy is needed to position the dots that are formed, and as the number of steps involved in varying the dot sizes is increased, greater dot size accuracy is required.\nThus, when a plurality of recording elements are employed, dot positioning errors and the use of non-uniform dot sizes can result in the deterioration of the image quality.\nIt is apparent that the demand for increased image quality can not be satisfied merely by improving the accuracy of the positioning of a carriage and a recording head and the accuracy in the production of dot sizes, so that accordingly, the shortcomings attributable to inaccurate dot positioning and to the unstable production of accurately sized dots are not resolved.\nIt is, therefore, one object of the present invention to provide at a low manufacturing cost an ink jet recording apparatus that can not only accurately detect the position of a recording head but can also accurately stabilize the positioning and the sizing of dots, a recording head therefor, and an element substrate to be used for the recording head.\nTo achieve the above object, according to one aspect of the present invention, provided is a recording head substrate on which are mounted energy generating elements that contribute to the formation of images by a recording head, and on which both light-receiving elements and light-emitting elements, or at least, light-receiving elements are mounted.\nThe light-receiving elements can be photodiodes or CCDs.\nIn addition, a controller for controlling the energy generating elements and the light-receiving elements is also mounted on the recording head substrate.\nIn this case, it is preferable that the light-receiving elements and at least one part of the controller be produced during the same manufacturing process.\nThe energy generating elements and the light-receiving elements are arranged along at least one line on the recording head substrate.\nThe energy generating elements and the light-receiving elements are arranged along a plurality of lines, and the lines are parallel to each other.\nIn this case, on the individual lines the number of the energy generating elements may be equal to the number of the light-receiving elements, but it is preferable that the number of the light-receiving elements be greater than the number of the energy generating elements.\nAccording to one more aspect of the present invention, provided is a recording head comprising:\nthe above described recording head substrate;\na top board in which are formed liquid flow paths that correspond to the energy generating elements; and\ndischarge orifices (port) which is communicated with the liquid flow path of the top plate and through which liquid is discharged by the application of energy by the energy generating elements,\nwherein the light-receiving elements and the light-emitting elements on the recording head substrate are optically opposite a face on which an image is formed by using the discharge ports.\nAccording to the present invention, as is described above the energy generating elements and the light-receiving elements are mounted on the same substrate. Therefore, when the light-receiving elements optically detect dots formed by the energy generating elements, accurate information concerning the positioning, the sizes and the densities of the image dots can be obtained quickly. Further, since in contrast to an arrangement where the energy generating elements, the light-emitting elements and the light-receiving elements are mounted on separate substrates, the process for the formation of the individual elements can be commonly employed and no connections are required, the manufacturing cost and the size of an apparatus can be considerably reduced.\nAccording to another aspect of the present invention, provided is a recording head substrate on which are mounted energy generating elements that contribute to the formation of images by a recording head, and on which are mounted a plurality of head position detecting elements for detecting the position of the recording head.\nAccording to an additional aspect of the present invention, a recording head, for forming images using energy generating elements, comprises:\na substrate on which are mounted not only the energy generating elements but also a head position detecting element for detecting the position of the recording head.\nAccording to a further aspect of the present invention, a liquid recording apparatus comprises:\na recording head for forming images employing energy generating elements while moving on a line;\nhead position detecting elements that are provided for the recording head for detecting the position of the recording head; and\na member in the recording apparatus that, in order to be detected by the head position detecting elements, is fixed opposite the head detecting element and along a track where the recording head moves.\nThe head position detecting elements are mounted on a substrate on which the energy generating elements are also mounted. In addition, it is preferable that, in accordance with position data for the recording head, detected by the head position detecting elements, and other recorded data, a circuit for generating signals to drive the energy generating elements, and light-receiving elements, for detecting an image that is formed, be mounted on the substrate on which the energy generating elements are mounted.\nThe head position detecting elements may be magnetic detecting elements, light-receiving elements or electric field detecting elements. The energy generating elements may be electro-thermal converting elements for heating liquid and inducing film boiling in order to discharge liquid droplets for forming images.\nAs is described above, according to the present invention, since the energy generating elements that contribute to image recording and the elements for detecting the position of the recording head are mounted on the same substrate, the accuracy at which the position of an image can be recorded is extremely high. In addition, since using semiconductor fabrication processing at least the elements having two functions can be mounted on the same substrate at the same time, the manufacturing costs can be drastically reduced."}
-{"text": "Receivers that detect stereophonic/monophonic signals are incorporated into a vast number of devices used in everyday life. For example, such receivers are used in automobile radios, a variety of communication systems such as cellular telephones, and even in children's toys. Unfortunately, many modern receiver systems suffer from performance shortfalls, such as frequent switchover between monophonic and stereophonic modes due to noisy channel environments and false detection of stereophonic signals as monophonic due to rated maximum system deviation (RMSD) mismatch.\nIn order to receive FM audio signals, be they music or any other type of information, a receiver must be robust enough to handle changes in the channel wherein the transmission could become very noisy and/or must overcome interference. Generally, a pilot tone is transmitted as part of the baseband signal that is used to modulate an FM carrier signal in order to indicate the nature of the transmission to be stereophonic. The energy of the pilot tone may fluctuate significantly in a harsh channel scenario. Simply comparing the pilot tone energy, estimated at the receiver, against a predetermined threshold may cause the receiver to switch between monophonic and stereophonic mode too frequently and degrade the entertainment quality of the audio program delivered to the consumer.\nIn addition, the receiver structure and the accompanying algorithms must also be flexible enough to handle a situation where the transmitted FM signal RMSD is not known beforehand. Generally, the allowed RMSD values are 75 kHz and 50 kHz. Hence, a mono/stereo transmission may be utilizing either one of them. However, the receiver may be set to operate at a different RMSD than the received signal RMSD. If the received signal RMSD and the RMSD of the receiver are not matched, a situation may occur wherein a stereophonic signal may be falsely interpreted as monophonic by the receiver. This results in the listener being denied the stereophonic quality of the program that the service provider is transmitting on the airwaves."}
-{"text": "As a measure against global warming, solar photovoltaic power generation has become popular in the world. For the solar photovoltaic power generation, a photoelectric conversion device (also called a solar cell) in which light energy is converted into electrical energy by using photoelectric characteristics of a semiconductor is applied in many cases, as compared to the case of utilizing solar heat.\nProduction of photoelectric conversion devices tends to increase year after year. For example, the total production of solar cells in the world in 2005 was 1,759 MW, which is a drastic increase of 147% as compared to that in the previous fiscal year. Photoelectric conversion devices which have become popular worldwide use a crystalline semiconductor; photoelectric conversion devices using a single crystal silicon substrate or a polycrystalline silicon substrate constitute a large part of the production.\nThe thickness of a crystal-type photoelectric conversion device using silicon, which is needed to absorb sunlight, is only about 10 \u03bcm. Nevertheless, a single crystal silicon substrate or a polycrystalline silicon substrate manufactured as a product has a thickness of about 200 to 300 \u03bcm. That is, the photoelectric conversion device using a single crystal semiconductor substrate or a polycrystalline semiconductor substrate has a thickness ten or more times as large as the thickness needed for photoelectric conversion, and thus the single crystal silicon substrate or the polycrystalline silicon substrate is far from being effectively utilized wholly. Speaking of extremes, most part of the single crystal silicon substrate or the polycrystalline silicon substrate functions only as a structural body for keeping the shape of the photoelectric conversion device.\nWith the increase of production of photoelectric conversion devices year after year, short of supply of polycrystalline silicon, which is the material of a silicon substrate, and resulting rise of cost of the same have become problems of the industry. The production of polycrystalline silicon is expected to be about 36,000 tons in 2007; in contrast, 25,000 tons or more of polycrystalline silicon is needed for semiconductor (LSI) and 20,000 tons or more of polycrystalline silicon is needed for solar cells, which means polycrystalline silicon seems to be in short of supply by about 10,000 tons. This short of supply is predicted to continue.\nThere are a variety of structures of photoelectric conversion devices. In addition to a photoelectric conversion device having a typical structure in which an n-type or a p-type diffusion layer is formed in a single crystal silicon substrate or a polycrystalline silicon substrate, a stacked photoelectric conversion device in which different kinds of unit cells, i.e., a unit cell formed of a single crystal semiconductor and a unit cell formed of an amorphous semiconductor, are combined is known (see Examined Patent Application Publication No. H6-044638). The photoelectric conversion devices are the same in that a single crystal silicon substrate or a polycrystalline silicon substrate is used. Here, as another mode of a photovoltaic device using a single crystal semiconductor substrate, a photovoltaic device using a single crystal semiconductor layer formed into a slice is given. For example, Patent Document 4 (Patent Document 4: Japanese Published Patent Application No. H10-335683) discloses a tandem solar cell in which hydrogen ions are implanted into a single crystal silicon substrate, a single crystal silicon layer which is separated from the single crystal silicon substrate in a layer shape is disposed over a support substrate in order to lower the cost and save resources while maintaining high conversion efficiency. In this tandem solar cell, a single crystal semiconductor layer and a substrate are bonded to each other with a conductive paste.\nOn the other hand, photoelectric conversion devices using a crystalline silicon thin film have also been developed. For example, a method for manufacturing a silicon thin-film solar cell in which a crystalline silicon film is deposited over a substrate by a plasma CVD method using a VHF of 27 MHz or higher which has been pulse modulated is described (see Japanese Published Patent Application No. 2005-50905). Further, a technique for controlling plasma process condition to optimize dopant concentration in crystal grains and crystal grain boundaries when a polycrystalline silicon thin film is formed by a plasma CVD method over a special electrode called a texture electrode which has minute unevenness on its surface is disclosed (see Japanese Published Patent Application No. 2004-14958). However, a crystalline thin-film silicon solar cell is still inferior to a single crystal silicon solar cell in crystal quality and photoelectric conversion characteristic. Moreover, a crystalline silicon film needs to be deposited to a thickness of 1 \u03bcm or more by a CVD method, which leads to a problem of low productivity."}
-{"text": "1. Field of the Invention\nThis invention relates to laser technology, and more particularly, to a polarization-type laser detection system which is capable of distinguishing between a reflected light signal from a target object and a back-scattered light signal from a mass of suspended particles in the air.\n2. Description of Related Art\nFIG. 1 is a schematic diagram of a conventional laser detection system. As shown, this laser detection system includes a control unit 10, a laser emitter 20, and an optical receiver 30. This laser detection system is used to detect whether any target object is present nearby. In the case of a target object 40 is present nearby, the laser beam from the laser emitter 20 will be reflected back by the target object 40 and received by the optical receiver 30. The received light is then analyzed by the control unit 10 to indicate the presence of the target object 40.\nOne drawback to the foregoing laser detection system, however, is that, under bad weather conditions when the air is filled with suspended particles, the emitted laser beam from the laser emitter 20 would be scattered back, causing the optical receiver 30 to received a back-scattered light signal that would make the control unit 10 unable to perform the intended target detection. The laser detection system would therefore operate improperly under bad weather conditions."}
-{"text": "1. Field of the Invention\nThis invention relates to a pressure sensor and particularly to an inexpensive and simple pressure sensor having good static pressure characteristics.\n2. Description of the Related Art\nA conventional pressure sensor has a silicon base 9 and a boss 10 (see, for example, JP-UM-A-5-50335). The conventional pressure sensor will now be described in detail with reference to FIG. 1. FIG. 1 is a sectional view of the conventional pressure sensor.\nIn FIG. 1, a metal 5, which is a metal base, has a hole 6 to which a pressure is applied. Also a glass 4, which is a base, has the hole 6 to which a pressure is applied. The glass 4 and the metal 5 are mounted to each other by Au eutectic bond or the like. The glass 4 performs electrical and mechanical insulation.\nAlso, the silicon base 9 has the hole 6 to which a pressure is applied. The silicon base 9 and the glass 4 are mounted to each other. The boss 10 is provided between the silicon base 9 and the glass 4.\nA silicon sensor 1, which is a sensor, has a diaphragm 2 connected to the hole 6. The silicon sensor 1 has a strain gauge 3 for converting a strain (displacement) occurring in the diaphragm 2 to an electric signal. One side of the silicon sensor 1 is mounted to the silicon base 9.\nThe other side of the silicon sensor 1 contacts a room B. The hole 6 and the diaphragm 2 form a room A.\nA static pressure is applied to the silicon sensor 1, the silicon base 9, the glass 4 and the other parts in the conventional example of FIG. 1.\nIn the conventional example of FIG. 1 constructed as described above, the pressure applied to the hole 6 is converted to an electric signal. The strain gauge 3 in the conventional example of FIG. 1 generates an electric signal based on the differential pressure between the room A and the room B and the static pressure.\nWhen the static pressure is applied, the silicon sensor 1, the silicon base 9 and the glass 4 are deformed, respectively. Since the silicon sensor 1 and the silicon base 9 have a large Young's modulus, these are deformed slightly. Since the glass 4 has a small Young's modulus, it is deformed largely.\nThe silicon base 9 and the boss 10 in the conventional example of FIG. 1 restrain transfer of the influence of deformation of the glass 4 to the strain gauge 3 of the silicon sensor 1. The boss 10 reduces the bond area between the silicon base 9 and the glass 4.\nMoreover, some conventional pressure sensors (semiconductor pressure converting apparatus) separately have a structure for detecting a differential pressure and a structure for detecting a static pressure, and also have a structure for reducing interference with a differential pressure signal while increasing output of a static pressure signal (see, for example, Japanese Patent No. 2,656,566).\nHowever, the conventional example of FIG. 1 has a problem that an error occurs in the application of the static pressure (static pressure characteristics).\nSpecifically, since the Young's modulus of silicon is different from the Young's modulus of glass, when the static pressure is applied, the deformation of the silicon sensor 1 and the silicon base 9 differs from the deformation of the glass 4, and a strain based on the deformation of the glass is generated in the strain gauge 3.\nMoreover, since the glass 4 has characteristics such as delayed elasticity and viscoelasticity, it causes a strain in the diaphragm 2 and thus causes a strain in the strain gauge 3. This causes an error in the conventional example of FIG. 1.\nAlso, the silicon base 9 and the boss 10 in the conventional example of FIG. 1 have problems of increase in the number of components, increase in the number of processing steps, and high cost.\nMoreover, the formation of the boss 10 has a problem of deteriorating the yield of bond. The formation of the boss 10 also has a problem of lowering broken pressure reduce the bond area.\nMeanwhile, the conventional example of Japanese Patent No. 2,656,566 has a problem that it does not restrain occurrence of an error in the static pressure characteristics and cannot acquire good static pressure characteristics."}
-{"text": "Butyric acid (BA) is a natural product. It is supplied to mammals from two main sources: 1) the diet, mainly from dairy fat, and 2) from the bacterial fermentation of unabsorbed carbohydrates in the colon, where it reaches mM concentrations (Cummings, Gut 22:763-779, 1982; Leder et al., Cell 5:319-322, 1975).\nBA has been known for nearly the last three decades to be a potent differentiating and antiproliferative agent in a wide spectra of neoplastic cells in vitro (Prasad, Life Sci. 27:1351-1358, 1980). In cancer cells, BA has been reported to induce cellular and biochemical changes, e.g., in cell morphology, enzyme activity, receptor expression and cell-surface antigens (Nordenberg et al., Exp. Cell Res. 162:77-85, 1986; Nordenberg et al., Br. J. Cancer 56:493-497, 1987; and Fishman et al., J. Biol. Chem. 254:4342-4344, 1979).\nAlthough BA or its sodium salt (sodium butyrate, SB) has been the subject of numerous studies, its mode of action is unclear. The most specific effect of butyric acid is inhibition of nuclear deacetylase(s), resulting in hyperacetylation of histones H3 and H4 (Riggs, et al., Nature 263:462-464, 1977). Increased histone acetylation following treatment with BA has been correlated with changes in transcriptional activity and the differentiated state of cells (Thorne et al., Eur. J. Biochem. 193:701-713, 1990). BA also exerts other nuclear actions, including modifications in the extent of phosphorylation (Boffa et al., J. Biol. Chem. 256:9612-9621, 1981) and methylation (Haan et al., Cancer Res. 46:713-716, 1986). Other cellular organelles, e.g., cytoskeleton and membrane composition and function, have been shown to be affected by BA (Bourgeade et al., J. Interferon Res. 1:323-332, 1981). Modulations in the expression of oncogenes and suppressor genes by BA were demonstrated in several cell types. Toscani et al., reported alterations in c-myc, p53 thymidine kinase, c-fos and AP2 in 3T3 fibroblasts (Oncogene Res. 3:223-238, 1988). A decrease in the expression of c-myc and H-ras oncogenes in B16 melanoma and in c-myc in HL-60 promyelocytic leukemia was also reported (Prasad et al., Biochem. Cell Biol. 68:1250-1255, 1992; and Rabizadeh et al., FEBS Lett. 328:225-229, 1993).\nBA has been reported to induce apoptosis, i.e., programmed cell death. SB has been shown to produce apoptosis in vitro in human colon carcinoma, leukemia and retinoblastoma cell lines (Bhatia et al., Cell Growth Diff. 6:937-944, 1995; Conway et al., Oncol. Res. 7:289-297, 1993; Hague et al.; Int. J. Cancer 60:400-406, 1995). Apoptosis is the physiological mechanism for the elimination of cells in a controlled and timely manner. Organisms maintain a delicate balance between cell proliferation and cell death, which when disrupted can tip the balance between cancer, in the case of over accumulation of cells, and degenerative diseases, in the case of premature cell losses. Hence, inhibition of apoptosis can contribute to tumor growth and promote progression of neoplastic conditions.\nThe promising in vitro antitumor effects of BA and BA salts led to the initiation of clinical trials for the treatment of cancer patients with observed minimal or transient efficacy. [Novogrodsky et al., Cancer 51:9-14, 1983; Rephaeli et al., Intl. J. Oncol. 4:1387-1391, 1994; Miller et al., Eur. J. Cancer Clin. Oncol. 23:1283-1287, 1987].\nClinical trials have been conducted for the treatment of .beta.-globin disorders (e.g., .beta.-thalassemia and sickle-cell anemia) using BA salts. The BA salts elevated expression of fetal hemoglobin (HbF), normally repressed in adults, and favorably modified the disease symptoms in these patients (Stamatoyannopouos et al., Ann. Rev. Med. 43:497-521, 1992). In this regard, arginine butyrate (AB) has been used in clinical trials with moderate efficacy (Perrine et al., N. Eng. J. Med. 328:81-86, 1993; Sher et al., N. Eng. J. Med. 332:1606-1610, 1995). The reported side effects of AB included hypokalemia, headache, nausea and vomiting in .beta.-thalassemia and sickle-cell anemia patients.\nButyric acid derivatives with antitumor activity and immunomodulatory properties have been reported in U.S. Pat. No. 5,200,553 and by Nudelman et al., 1992, J. Med. Chem. 35:687-694. The most active butyric acid prodrug reported in these references was pivaloyloxymethyl butyrate (AN-9). None of the compounds disclosed in these references included carboxylic acid-containing oxyalkyl compounds of this invention.\nBA and/or its analogues have also been reported to increase the expression of transfected DNA (Carstea et al., 1993, Biophys. Biohem. Res. Comm. 192:649; Cheng et al., 1995, Am. J. Physical 268:L615-L624) and to induce tumor-restricted gene expression by adenovirus vectors (Tang et al., 1994, Cancer Gene Therapy 1:15-20). Tributyrin has been reported to enhance the expression of a reporter gene in primary and immortalized cell lines (Smith et al., 1995, Biotechniques 18:852-835).\nHowever, BA and its salts are normally metabolized rapidly and have very short half-lives in vivo, thus the achievement and maintenance of effective plasma concentrations are problems associated with BA and BA salts, particularly for in vivo uses. BA and BA salts have required large doses to achieve even minimal therapeutic effects. Because of the high dosage, fluid overload and mild alkalosis may occur. Patients receiving BA emanate an unpleasant odor that is socially unacceptable.\nWhile BA salts have been shown to increase HbF expression, and appear to hold therapeutic promise with low toxicity in cancer patients, they nevertheless have shown low potency in in vitro assays and clinical trials. There also remains a need to identify compounds as effective or more effective than BA or BA salts as differentiating or anti-proliferating agents for the treatment of cancers. Such compounds need to have higher potency than BA without the problems associated with BA (such as bad odor). Consequently, there remains a need for therapeutic compounds that either deliver BA to cells in a longer acting form or which have similar activity as BA but a longer duration of effectiveness in vivo.\nThe compounds of this invention address these needs and are more potent than BA or BA salts for treating cancers and other proliferative diseases, for treating gastrointestinal disorders, for wound healing and for treating blood disorders such as thalassemia, sickle cell anemia and other anemias, for modulating an immune response, for enhancing recombinant gene expression, for treating insulin-dependent patients, for treating cystic fibrosis patients, for inhibiting telomerase activity, for detecting cancerous or malignant cells, for treating virus-associated tumors, especially EBV-associated tumors, for augmenting expression of a tumor suppressor gene and for inducing tolerance to an antigen. For example, one of the advantages of the compounds of the invention is increased water solubility of the free carboxylic acids compounds of the invention and their salts, and easier administration, especially for intravenous administration."}
-{"text": "1. Field\nOne or more embodiments disclosed herein relate to a scan driving circuit and a driving method for the scan driving circuit.\n2. Description of the Related Art\nPersonal computers, portable phones, portable information terminals, and monitors of various information devices are all equipped with displays. The display may be, for example, a liquid crystal display, an organic light emitting diode display, or a plasma display panel. An organic light emitting diode display in particular has excellent emission efficiency, luminance, viewing angle, and response time.\nAn organic light emitting diode display generates images using pixels that include organic light emitting diodes. The pixels are arranged in a matrix at cross points of data lines, scan lines, and power supply lines. In addition to an organic light emitting diode, each pixel includes a driving transistor and one or more capacitors. Light is emitted from the diode based on a recombination of electrons and holes in a light emitting active layer.\nOver time, the amount of current flowing in the organic light emitting diode of each pixel may change based on a deviation in the threshold voltage of the driving transistor. As a result, non-uniformity may occur in the display. Additionally, the characteristics of the driving transistor may change based on manufacturing process parameters. Because it is difficult for the transistors in an organic light emitting diode display to have the same characteristics, the threshold voltage deviation of the driving transistors in the pixels may be different.\nIn attempt to overcome this deviation, a compensation circuit may be used in each pixel. The compensation circuit may charge a voltage corresponding to the threshold voltage of the driving transistor for 1 horizontal period. However, in some displays and especially in a high-resolution organic light emitting diode display, coupling may occur in a parasitic capacitor between lines and horizontal striped patterns. This may adversely affect display performance."}
-{"text": "Conventionally, as disclosed in Patent Literature 1, an on-board electronic device having CPUs and substrates, respectively corresponding to types of functions, such as a navigation board on which a CPU to perform navigation processing is mounted and an audio board on which a CPU to perform audio processing is mounted, is known. Note that this type of on-board electronic device is often integrated with a display and a touch panel.\nFurther, in recent years, in accordance with multimedization in the vehicle, a configuration, where a user interface unit (UI unit) having integrated display and touch panel is connected to a plurality of electronic devices mounted on the vehicle such as a navigational device and an audio instrument, via a display control unit (hereinbelow, DCU), is increasingly adopted. The DCU plays a role in control of the respective operations of the plurality of electronic devices connected to the own device, and generation (or acquisition) of image data to be displayed on the display based on requests from the respective electronic devices, and display of the image data on the display. Hereinbelow, a unit where the UI unit and the DCU are integrated will be referred to as a display unit for vehicle.\nA conventional DCU has two substrates, i.e., a main board mounted with a main CPU and a sub board mounted with a sub CPU and a power supply circuit. The main CPU mainly controls the operations of the electronic devices connected to the own device, and draws an image to be outputted to the display. The sub CPU performs management (acquisition and storage) of vehicle information inputted from an in-vehicle network, and controls electric power supply to the respective elements of the display unit for vehicle in cooperation with the power supply circuit.\nThe conventional DCU has two substrates, i.e. the main board and the sub board, mainly for reducing the area of each substrate to an area smaller than that of the display.\nMore particularly, to realize a function, performed with a main board and a sub board, using one substrate, the area of the substrate is larger than the area of the display. When the area of the substrate is larger than the area of the display, its housing is also larger than the display, to impair mountability to the vehicle.\nAccordingly, in the conventional DCU, a function to be provided in the DCU is shared with the two substrates, the main board and the sub board, so that the size of each substrate is suppressed to be equal to or smaller than the size of the display.\nHowever, in the conventional display unit for vehicle (more particularly, DCU) having two substrates, the manufacturing cost such as costs of parts and machining cost is increased. The machining cost here includes costs necessary for assembling work of the respective substrates in the housing and wiring connection work between the substrates.\nFurther, it is desired to increase the CPU clock for improvement in DCU performance. However, in the conventional DCU, the space in the housing is partitioned with the two substrates. When the clock is increased, radiation with a cooling fan becomes insufficient. Accordingly, from the viewpoints of manufacturing cost and heat dissipation, it is preferable that the number of the substrates is one.\nOn the other hand, as described above, when a function realized with two substrates is simply integrated in one substrate, the area of the substrate becomes larger than the area of the display, which might impair the mountability. Especially, when the difference between the area of the display and the area of the substrate is large, i.e., when the area of the substrate is larger, the mountability is impaired."}
-{"text": "Network nodes forward data. Network nodes may take form in one or more routers, one or more bridges, one or more switches, one or more servers, or any other suitable communications processing device. The data is commonly formatted as packets and forwarded using forwarding tables. A packet is a formatted unit of data that typically contains control information and payload data. Control information may include: information that identifies sources and destinations, such as addresses, error detection codes like checksums, sequencing information, etc. Control information is typically found in packet headers and trailers. Payload data is typically located between the packet headers and trailers.\nForwarding packets involves various processes that, while simple in concept, can be complex. The processes involved in forwarding packets vary, depending on the type of forwarding method used. Multicast is the preferred method of data forwarding for many networks. One reason for this is that multicast is a bandwidth-conserving technology that reduces traffic by simultaneously delivering data to multiple receivers. However, some network environments are not well suited to support multicast. Doing so in such environments often involves discovering and maintaining significant amounts of control, or state, information. Setting up and maintaining this control information has a tendency to become complex and costly in terms of computing resources, and can become a major limiting factor in overall network performance. Another issue with multicast is that due to packet delivery mechanisms used, packets are sometimes forwarded to locations where the packets were not desired. This unnecessary delivery of packets represents an unwelcome burden on network performance. Overcoming this burden by traditional means involves generation and maintenance of even more control information."}
-{"text": "Amplitude-shift-keyed modulation is a commonly known technique used to modulate rf and optical carriers for the transmission of data. A typical transmission involves sending bursts of carrier signals, the presence of a burst identifying a binary one or zero and the absence of a burst or carrier being representative of the opposite binary digit. Examples of patents describing various ASK modulations are U.S. Pat. Nos. 4,947,407, 4,860,320, and 4,829,560.\nTechniques for improving detection of data signals detected with optical receivers are known in the art. In one such technique as shown and described in U.S. Pat. No. 4,431,916 delayed versions of input signals are compared with undelayed versions to reproduce data signals. Threshold values are created from these signals. U.S. Pat. No. 4,507,795 describes an apparatus for locating leading and trailing edges of pulses derived from an RF input.\nOne problem associated with ASK receivers involves the large amplitude swings of incident RF or optical carrier signals. Such large swings tend to introduce undesirable variations in the output data pulse width. Such variations often are due to the receiver's bandpass filtering which causes sloping, rising and falling edges at the output of the logarithmic amplifier that is commonly used in receivers."}
-{"text": "This invention relates to bioadhesive compositions, particularly electrically conductive hydrogel compositions having bioadhesive properties. The invention further relates to biomedical skin electrodes incorporating such hydrogel bioadhesive compositions that are electrically conductive.\nBiomedical skin electrodes are widely used in a variety of situations, whenever for example it is required to establish an electrical connection between the surface of the body of the patient and external medical equipment for transmission of electrical signals.\nModem medicine uses many medical procedures where electrical signals or currents are received from or delivered to a patient\"\"s body. The interface between medical equipment used in these procedures and the skin of the patient is usually some sort of biomedical electrode. Such electrodes typically include a conductor which must be connected electrically to the equipment, and a conductive medium adhered to or otherwise contacting skin of the patient, and they are of varying types with a wide variety of design configurations which will generally depend on their intended use and whether for example they are to be used as transmission electrodes or sensing i.e. monitoring electrodes.\nAmong the therapeutic procedures using biomedical electrodes are transcutaneous electric nerve stimulation (TENS) devices used for pain management; neuromuscular stimulation (NMS) used for treating conditions such as scoliosis; defibrillation electrodes to dispense electrical energy to a chest cavity of a mammalian patient to defibrillate heart beats of the patient; and dispersive electrodes to receive electrical energy dispensed into an incision made during electrosurgery.\nAmong diagnostic procedures using biomedical electrodes are monitors of electrical output from body functions, such as electrocardiograms (ECG) for monitoring heart activity and for diagnosing heart abnormalities.\nFor each diagnostic, therapeutic, or electrosurgical procedure, at least one biomedical electrode having an ionically conductive medium containing an electrolyte is adhered to or is otherwise contacted with mammalian skin at a location of interest and is also electrically connected to electrical diagnostic, therapeutic, or electrosurgical equipment. A critical component of the biomedical electrode is the conductive medium which serves as the interface between the mammalian skin and the diagnostic, therapeutic, or electrosurgical equipment, and which is usually an ionically conductive medium.\nBiomedical electrodes are used among other purposes to monitor and diagnose a patient\"\"s cardiovascular activity. Diagnostic electrodes are used to monitor the patient immediately and are only applied to the patient for about five to ten minutes. Monitoring electrodes, however, are used on patients in intensive care for up to three days continuously. In contrast, Holter electrodes are used to monitor a patient during strenuous and daily activities.\nAlthough all of the biomedical electrodes just referred to are used to record cardiovascular activity, each electrode requires specific features or characteristics to be successful. Thus, the diagnostic electrode does not have to remain adhered to a patient for extensive periods but it does have to adhere to hairy, oily, dry and wet skin effectively for the five to ten minutes of use. The monitoring electrode has to adhere for a longer period of time although the patient is often immobile during the monitoring period. The Holter electrodes is susceptible to disruption from adhesion due to physical motion, perspiration, water, etc., and therefore requires the best adhesion and at the same time comfort and electrical performance.\nIn the biomedical electrodes known in the prior art the ionically conductive medium which serves as an interface, between the skin of a mammalian patient and the electrical instrumentation, ranges from conductive gels and creams to conductive pressure sensitive adhesives. However, while the conductive media can be in the form of pressure sensitive conductive adhesives, for monitoring or Holter biomedical electrode use such conductive adhesives are not generally adequate on their own to maintain adhesion to mammalian skin and additional hypoallergenic and hydrophobic pressure sensitive adhesives may be employed around the conductive medium to provide the required mammalian skin adhesion. U.S. Pat. No. 5,012,810 (Strand et al.) and U.S. Pat. Nos. 4,527,087, 4,539,996, 4,554,924 and 4,848,353 (all Engel) are examples of documents that disclose biomedical electrodes which have a hydrophobic pressure sensitive adhesive surrounding the conductive medium.\nIn general, a desirable skin electrode is one which maintains good electrical contact with the skin and is free of localised current hot spots, i.e. exhibits uniform conductivity. For example, it has been found that a prior art electrode utilising karaya gum tends to creep in use and flatten out, exposing skin to possible direct contact with the current distribution member or lead wire. A desirable skin electrode should also usually have a low electrical impedance.\nIt is an object of this invention to provide hydrogel adhesives possessing controlled and predictable adhesive properties which may be readily varied to suit different uses and, in the case of medical electrodes or similar devices, different configurations or applications. It is also an object of the invention to provide such hydrogel adhesives which in addition may possess superior electrical characteristics as compared to those commonly associated with bioadhesive hydrogels."}
-{"text": "Efficient list decoding beyond half a minimum distance for Reed-Solomon and Bose, Ray-Chaudhuri and Hocquenghem (i.e., BCH) codes were first devised in 1997 and later improved almost three decades after the inauguration of an efficient hard-decision decoding method. In particular, for a given Reed-Solomon code C(n,k,d), a Guruswami-Sudan decoding method corrects up to n\u2212\u221a{square root over (n(n\u2212d))} errors, which effectively achieves a Johnson bound, a general lower bound on the number of errors to be corrected under a polynomial time for any code. Schmidt, Sidorenko, and Bossert devised a multi-sequence shift-register synthesis to find an error locator polynomial beyond half a minimum distance for low-rate (i.e., <\u2153) Reed-Solomon codes when a unique solution exists. Apart from small probability of failure due to ambiguity, the resulting decoding radius is identical to that of the Sudan technique. The Sudan technique extended the Guruswami-Sudan technique to achieve subfield Johnson bounds for subfield subcodes of Reed-Solomon codes by distributing multiplicities across the entire subfield.\nWu presented a list decoding technique for Reed-Solomon and binary BCH codes. The Wu list decoding technique casts a list decoding as a rational curve fitting problem utilizing polynomials constructed by a Berlekamp-Massey technique. The Wu technique achieves the Johnson bound for both Reed-Solomon and binary BCH codes. Beelen and Hoeholdt re-interpreted the Wu list decoding technique in terms of Gr\u00f6bner bases and an extended Euclidean technique to list decoded binary Goppa codes up to the binary Johnson bound.\nIt would be desirable to implement a combined Wu and Chase decoding of cyclic codes."}
-{"text": "1. Field of the Invention\nThe systems and methods of this invention generally related to communications systems. In particular, the systems and methods of this invention relate to providing a variable state length initialization.\n2. Description of Related Art\nMulticarrier modulation, which is also known as Discrete Multitone Transmission (DMT), transceivers step a through a number of initialization states prior to entering steady-state communication or \u201cshowtime.\u201d In particular, these various initialization states include channel discovery, transceiver training, channel analysis, and the like. These various initialization states allow, for example, the determination of transmitter power levels, line characteristics, training of receiver function such as equalizers or echo cancellers, or any other feature necessary to establish communication, or to exchange parameters and settings, between transceivers."}
-{"text": "The present invention relates to a teletext converter which is external to a television receiver (\"set-top\" converter), and more particularly, to such converters having a \"transparent mode\".\nTeletext is a service that broadcasts digital information within the vertical blanking interval (VBI) of a standard broadcast television signal and at the receiver presents the information as digitally generated text and/or pictures on a display screen. In the most commonly considered mode of operation, the standard TV picture is completely replaced by the digitally generated picture. \"Transparent mode\" refers to the situation in which some region of the screen, rather than being defined by the teletext signal, is occupied by the regular video signal. One example is a regular video picture having a teletext generated sub-title or caption. Ideally, the teletext converter, which converts the digital signals into video signals, is located within the receiver. However, almost all existing receivers do not have such an internal converter. Thus an external set-top converter is used if teletext signal reception is desired on most existing sets. A set-top converter receives the antenna signal, converts the digital teletext signal into a video signal, and then provides an output signal for connection to the receivers antenna input terminal.\nOne prior art method of producing a transparent mode set-top converter is to decode the regular video to baseband red, green and blue signals, combine these signals with R, G, and B signals derived from the decoded teletext, re-encode to composite NTSC, and then modulate with the composite NTSC signal an output carrier having a selected channel frequency. This approach suffers the deficiencies of high cost and significant signal quality degradation.\nIn particular, the decoding of composite video to baseband color signals, involving the separation of luminance and chrominance, band-limiting these signals and various other distortion producing processes degrades the fidelity of the standard TV picture produced after re-encoding and finally decoding again in the TV receiver.\nIt is therefore desirable to provide a set-top teletext decoder with transparent mode that does not cause signal degradation and has accurate color fidelity."}
-{"text": "The present invention relates to an ink-jet print method and apparatus for printing lines by discharging an ink from an ink-jet head onto a recording member, a color filter, a display device, an apparatus having the display device, an ink-jet head unit adjusting device, and an ink-jet head unit.\nWith recent advances in personal computers, especially portable personal computers, the demand tends to arise for liquid crystal displays, especially color liquid crystal displays. However, in order to further popularize the use of liquid crystal displays, a reduction in cost must be achieved. Especially, it is required to reduce the cost of a color filter which occupies a large proportion of the total cost. Various methods have been tried to satisfy the required characteristics of color filters while meeting the above requirements. However, any method capable of satisfying all the requirements has not been established. The respective methods will be described below.\nThe first method is a pigment dispersion method. In this method, a pigment-dispersed photosensitive resin layer is formed on a substrate and patterned into a single-color pattern. This process is repeated three times to obtain R, G, and B color filter layers.\nThe second method is a dyeing method. In the dyeing method, a water-soluble polymer material as a dyeable material is applied onto a glass substrate, and the coating is patterned into a desired shape by a photolithographic process. The obtained pattern is dipped in a dye bath to obtain a colored pattern. This process is repeated three times to form R, G, and B color filter layers.\nThe third method is an electrodeposition method. In this method, a transparent electrode is patterned on a substrate, and the resultant structure is dipped in an electrodeposition coating fluid containing a pigment, a resin, an electrolyte, and the like to be colored in the first color by electrodeposition. This process is repeated three times to form R, G, and B color filter layers. Finally, these layers are calcined.\nThe fourth method is a print method. In this method, a pigment is dispersed in a thermosetting resin, a print operation is performed three times to form R, G, and B coatings separately, and the resins are thermoset, thereby forming colored layers. In either of the above methods, a protective layer is generally formed on the colored layers.\nThe point common to these methods is that the same process must be repeated three times to obtain layers colored in three colors, i.e., R, G, and B. This causes an increase in cost. In addition, as the number of processes increases, the yield decreases. In the electrodeposition method, limitations are imposed on pattern shapes which can be formed. For this reason, with the existing techniques, it is difficult to apply this method to TFTs. In the print method, a pattern with a fine pitch is difficult to form because of poor resolution and poor evenness.\nIn order to eliminate these drawbacks, methods of manufacturing color filters by an ink-jet system are disclosed in Japanese Patent Laid-Open Nos. 59-75205, 63-235901, and 1-217320. In these methods, inks containing coloring agents of three colors, i.e., R (red), G (green), and B (blue), are sprayed on a transparent substrate by an ink-jet system, and the respective inks are dried to form colored image portions. In such an ink-jet system, R, G, and B pixels can be formed at once, allowing great simplification of the manufacturing process and a great reduction in cost.\nWhen a color filter is to be manufactured by such an ink-jet system, an ink is discharged onto each pixel while an elongated ink-jet head having a plurality of ink discharging nozzles is scanned over a color filter substrate. This scanning is performed a plurality of number of times to color the respective pixel portions. In this case, the amounts of ink discharged from the respective ink discharging nozzles slightly differ from each other. If, therefore, each pixel array is colored with the same nozzle, the pixel arrays colored with the nozzles from which the ink is discharged in large amounts become dense in color, but the pixel arrays colored with the nozzles from which the ink is discharged in small amounts become light in color. Consequently, the resultant color filter has color irregularity.\nIn addition, to manufacture a color filter, such elongated ink-jet heads must be prepared for three colors, i.e., R (red), G (green), and B (blue). It takes a lot of time and labor to adjust the relative positions of these three heads."}
-{"text": "1. Field of the Invention\nThe present invention relates to a chemically amplified resist composition, and more particularly, to a photosensitive polymer having a cyclic backbone, and a resist composition for an ArF excimer laser obtained therefrom.\n2. Description of the Related Art\nAs semiconductor devices become highly integrated and complicated to fabricate, fine pattern formation is required.\nFurther, as the capacity of a semiconductor device increases to exceed 1 giga bit, a pattern size having a design rule of less than 0.2, .mu.m is required. Accordingly, there are limitations in using a conventional resist material with a KrF excimer laser (248 nm). Thus, a new resist material capable of being developed using an ArF excimer laser (193 nm) has been developed in a lithography process.\nThe resist material used in the lithography process using the ArF excimer laser has several problems in being commercially used, compared to the conventional resist materials. The most typical problems are transmittance of a polymer and resistance to dry etching. As the widely known ArF resist materials, (meth)acrylate polymers are generally used. In particular, the most typical resist material is a poly(methyl methacrylate-tert-butyl methacrylate-methacrylic acid) terpolymer system manufactured by IBM, Inc. However, such polymers have very weak resistance to dry etching.\nAccordingly, to increase the resistance to dry etching, a polymer having a backbone composed of an alicyclic compound such as an isobornyl group, an adamantyl group or a tricyclodecanyl group, is used. However, the resulting resist still exhibits weak resistance to dry etching."}
-{"text": "A need exists for a high speed screening device that can be installed in a new and/or existing facility and efficiently screen materials to a tight specification at an industrial minerals processing facilities.\nA further need exists for a screening device that can make multiple gradation cuts, of particulate simultaneously while creating very little dust in the facility.\nA need exists for a device that will prevent human harm during the screening of particulate.\nThe present embodiments meet these needs.\nThe present embodiments are detailed below with reference to the listed Figures."}
-{"text": "As computers and data processing equipment have grown in capability, users have developed applications that place increasing demands on the equipment. Thus, there is a continually increasing need to process more information in a given amount of time. One way to process more information in a given amount of time is to process each element of information in a shorter amount of time. As that amount of time is shortened, it approaches the physical speed limits that govern the communication of electronic signals. While it would be ideal to be able to move electronic representations of information with no delay, such delay is unavoidable. In fact, not only is the delay unavoidable, but, since the amount of delay is a function of distance, the delay varies according to the relative locations of the devices in communication.\nSince there are limits to the capabilities of a single electronic device, it is often desirable to combine many devices, such as memory components, to function together to increase the overall capacity of a system. However, since the devices cannot all exist at the same point in space simultaneously, consideration must be given to operation of the system with the devices located diversely over some area.\nTraditionally, the timing of the devices' operation was not accelerated to the point where the variation of the location of the devices was problematic to their operation. However, as performance demands have increased, traditional timing paradigms have imposed barriers to progress.\nOne example of an existing memory system uses DDR (double data rate) memory components. The memory system includes a memory controller and a memory module. A propagation delay occurs along an address bus between the memory controller and the memory module. Another propagation delay occurs along the data bus between the memory controller and the memory module.\nThe distribution of the control signals and a control clock signal in the memory module is subject to strict constraints. Typically, the control wires are routed so there is an equal length to each memory component. A \u201cstar\u201d or \u201cbinary tree\u201d topology is typically used, where each spoke of the star or each branch of the binary tree is of equal length. The intent is to eliminate any variation of the timing of the control signals and the control clock signal between different memory components of a memory module, but the balancing of the length of the wires to each memory component compromises system performance (some paths are longer than they need to be). Moreover, the need to route wires to provide equal lengths limits the number of memory components and complicates their connections.\nIn such DDR systems, a data strobe signal is used to control timing of both data read and data write operations. The data strobe signal is not a periodic timing signal, but is instead only asserted when data is being transferred. The timing signal for the control signals is a periodic clock. The data strobe signal for the write data is aligned to the clock for the control signals. The strobe for the read data is delayed by delay relative to the control clock equal to the propagation delay along the address bus plus the propagation delay along the data bus. A pause in signaling must be provided when a read transfer is followed by a write transfer to prevent interference along various signal lines used. Such a pause reduces system performance.\nSuch a system is constrained in several ways. First, because the control wires have a star topology or a binary tree routing, reflections occur at the stubs (at the ends of the spokes or branches). The reflections increase the settling time of the signals and limit the transfer bandwidth of the control wires. Consequently, the time interval during which a piece of information is driven on a control wire will be longer than the time it takes a signal wavefront to propagate from one end of the control wire to the other. Additionally, as more modules are added to the system, more wire stubs are added to each conductor of the data bus, thereby adding reflections from the stubs. This increases the settling time of the signals and further limits the transfer bandwidth of the data bus.\nAlso, because there is a constraint on the relationship between the propagation delays along the address bus and the data bus in this system, it is hard to increase the operating frequency without violating a timing parameter of the memory component. If a clock signal is independent of another clock signal, those clock signals and components to which they relate are considered to be in different clock domains. Within a memory component, the write data receiver is operating in a different clock domain from the rest of the logic of the memory component, and the domain crossing circuitry will only accommodate a limited amount of skew between these two domains. Increasing the signaling rate of data will reduce this skew parameter (when measured in time units) and increase the chance that a routing mismatch between data and control wires on the board will create a timing violation.\nAlso, most DDR systems have strict limits on how large the address bus and data bus propagation delays may be (in time units). These are limits imposed by the memory controller and the logic that is typically included for crossing from the controller's read data receiver clock domain into the clock domain used by the rest of the controller. There is also usually a limit (expressed in clock cycles) on how large the sum of these propagation delays can be. If the motherboard layout makes this sum too large (when measured in time units), the signal rate of the system may have to be lowered, thereby decreasing performance.\nIn another example of an existing memory system, the control wires and data bus are connected to a memory controller and are routed together past memory components on each memory module. One clock is used to control the timing of write data and control signals, while another clock is used to control the timing of read data. The two clocks are aligned at the memory controller. Unlike the previous prior art example, these two timing signals are carried on separate wires.\nIn such an alternate system, several sets of control wires and a data bus may be used to intercouple the memory controller to one or more of the memory components. The need for separate sets of control wires introduces additional cost and complexity, which is undesireable. Also, if a large capacity memory system is needed, the number of memory components on each data bus will be relatively large. This will tend to limit the maximum signal rate on the data bus, thereby limiting performance.\nThus, a technique is needed to coordinate memory operations among diversely-located memory components."}
-{"text": "1. Field of the Invention\nThis invention relates to airbrasive devices, and more particularly to devices for containing abrasive materials expelled by a gas-abrasive apparatus. The invention is particularly useful for dental applications.\n2. Background of the Invention\nThe use of sandblasting devices to contact various surfaces has been known for some time. These devices are also known in the art as airbrasive or air-abrasive devices. Such devices vary in size and design depending on the particularly utility desired.\nOne area where use of these devices has proved advantageous is in the etching or abrading of small surfaces. Devices designed for this use are typically hand held and capable of delivering fine streams of air-abrasive material through narrow nozzles.\nA number of decades ago, the use of air-abrasive devices gained favor in the dental art The methods developed were termed \"airbrasive techniques\" and were designed to supplement the use of traditional dental drills to prepare a tooth for cavity repair, prophylaxis or other methods that required that a portion of the tooth be removed or that required the roughing of a tooth surface. The advantage of using air-abrasive techniques is that the dental patient experiences less trauma to the oral cavity due to the absence of perceptible pressure, vibration, noises created by the contact of a drill to tooth enamel, and heat created by frictional forces. This has resulted in reduced pain, apprehension, and fear by patients.\nOne disadvantage of the use of air-abrasive dental apparatus is that abrasive materials are dispersed into the oral cavity during use in a relatively uncontrolled fashion, can be inhaled by the patient, and are difficult to remove after a procedure is complete. Another disadvantage is that such particles can be dispersed into the air and create a hygiene problem. Abrasive particles can carry pathogens and blood particles from the mouth and permit those pathogens and blood particles to contact otherwise uncontaminated surfaces.\nSomewhat similar disadvantages exist with use of air-abrasive devices in other applications. Often it is desirable to prevent abrasive materials from contacting surfaces proximate to the target surface, from accumulating abrasive material on the target surface area, or from permitting fine abrasive particles from becoming airborne.\nSeveral devices have been developed to affect the dispersion of abrasive particles within the oral cavity. Coston, U.S. Pat. No. 5,197,876 discloses a splatter guard for air polishing dental devices. The guard comprises a bell-shaped flexible cone that is attached to the end of an air-abrasive device and guides abrasive particles towards the surface being treated. Ho, U.S. Pat. No. 5,356,292 discloses a dental sandblasting confiner in the form of a flexible transparent cup. The nozzle of a sandblasting device can be inserted in large opening of the cup which forms a mold around the nozzle. The Ho device contains additional openings for access to a tooth surface and for discharging output. Lokken, U.S. Pat. No. 4,611,992 discloses an anti-splash device that can be attached to a dental tool. The device comprises an inverted U-shaped member with legs for attaching the device to the dental tool. Wright, U.S. Pat. No. 4,850,868 discloses a spray shield comprising a modified tube that can be attached to the end of a dental handpiece. The device is used to direct material dispensed form the handpiece in a controlled fashion so as to minimize the amount of airborne particles.\nWhile the above cited inventions address one or more of the described disadvantages of air-abrasive systems, they are subject to several detrimental limitations. Although minimizing the amount of abrasive material released, by guiding it downward for instance, has certain benefits, it is more preferable to contain a substantial portion of released abrasive material and permit facile removal. Many of the devices in the prior art guide, but do not completely contain abrasive material nor permit easy removal thereof. Other devices that do permit removal of abrasive material are obtrusive and interfere with visualization of the surface to be abraded, making it difficult to perform precise dental procedures. Furthermore, those devices that do permit removal of abrasive material typically rely on a vacuum source to remove that material. Such a vacuum source adds additional expense and can also be intrusive.\nThus, there is a need for a device that can contain a substantial portion of the abrasive material expelled from an air-abrasive device while not obstructing visualization of the surface to be abraded and permitting removal of the abrasive material without the aid of a vacuum source."}
-{"text": "The invention is based on a method and an apparatus as generally set forth hereinafter. Means for controlling the spring stiffness, for instance of a motor vehicle, are known (German patent No. 16 30 058); in this known apparatus, two work chambers of a shock absorber or telescoping spring are connected via external lines to an apparatus comprising a pump and two reservoirs. Only one-way check valves are disposed in the connecting lines to the telescoping spring. However, with this type of apparatus, the damper stiffness of this kind of shock absorber cannot be varied, because to do so energy must be supplied from outside--via the pump--which takes a relatively long time and means that there is a certain energy demand. Controlling the damper stiffness in a shock absorber is also known from German Offenglegungsschrift No. 33 04 815. Suspension systems of present-day vehicle types, in particular passenger vehicles, are typically optimized to an average operational case in terms of the spring stiffness and damper stiffness, with parameters being structurally fixed and remaining unchanged, except for effects associated with aging, during driving. Since in extreme operational cases, such as an empty or fully loaded vehicle, or with varying vehicle movement parameters (rapid cornering, braking, acceleration, smooth highway driving, and the like), optimal suspension or damping of the suspension system is not attained in all such operating situations, it is also already known to switch over among a plurality of damper settings or spring stiffnesses. This prevents long-term adaptation of the suspension system (and possibly the damper system, which is either arbitrarily integrated into the suspension system itself or is a component thereof), and especially it is impossible to make an automatic, finely-tuned adjustment to various road conditions or various kinds of driving, because a switchover in the characteristic curves of the suspension and/or damping system can be made only in stages, typically between only two operating states.\nIt is also known (U.S. Pat. No. 3,807,678), in a suspension system involving two masses, one of which may be one or more wheels of a vehicle and the other the vehicle body, to dispose a standard, passive compression spring between the two masses, which is called a passive isolating element and has a so-called active damper switched parallel to it. This damper, in which a piston slides in a cylinder and divides it into two work chambers, is considered to be active because an intervention is made into the damping properties, that is, into the positive volume displacements of the pressure medium in the various working halves of the damper by control means, in a so-called active manner. To this end, the two working chambers are connected crosswise and parallel to each other via opposed valves allowing a flow of pressure medium in only one direction; the amount of pressure medium then allowed to pass through these valves then also becomes \"active\" by appropriate control of the valves by means of suitably prepared sensor signals. Because in this known suspension system, the spring itself is entirely passive, but the damper is conceived of as being active in terms of its properties, the overall system in this patent is called a semiactive system. However, this term is not semantically related to the dampers of the present invention, which without reference to suspension systems not taken into account are themselves designated as so-called semiactive dampers, for reasons to be explained hereinafter.\nIt is also known, in wheel suspensions in vehicles, to provide so-called active damping means (see the article, \"Active Damping in Road Vehicle Suspension Systems\", published in the periodical, Vehicle System Dynamics, 12 (1983), pages 291-316). This publication is referred to also because it includes basic concepts, in theoretically detailed form, applicable in particular to active damping properties."}
-{"text": "Hair transplant procedures have been carried out for decades. Initially, a punch was used to remove a circular area of hairy skin containing ten or more follicular units (of 1-4 hairs each). The area of hairy skin replaced a like area of bald skin removed from the patient. Several of such \u201cplugs,\u201d were placed into areas in the bald part of the head.\nThe circular punch tool was later replaced by a hollow powered drill and the space left in the donor area was left to heal naturally. Both of these prior art procedures allowed wounds to stay open for weeks at a time exposing a patient to the discomfort from large wounds measuring 3-5 millimeters in diameter.\nToday there are two standard procedures for harvesting hair, the first involves a linear incision which permits the removal of a strip of hairy skin down into the fatty level of one quarter inch and measuring a number of square inches. The resulting wound is sutured closed and the strip is dissected into grafts (under a microscope), cooled in an ice bath or refrigerator and then transplanted into a bald area in needle size holes. Forceps grasp each graft and places them into holes in the bald area. In one form for dissection the hair from the strip of scalp uses blind harvesting of grafts from the strip of hairy skin which can result in significant damage to the hair. The damage occurs because the strip of hair in the hairy skin is forced through a cutting grid in order to make grafts of a predetermined size. The cutting blades of the grid are positioned at the most ideal distance between follicles. Unfortunately, the distance between follicles varies randomly. The result is that a significant number of the hair follicles can be damaged and die.\nThe second harvesting technique involves the use of a microscope to dissect the hair from the excised scalp.\nA third harvesting technique uses a punch which cores out from the scalp, the basic anatomical unit of hair, the follicular unit, which contains between one and four hairs each. The problems with this technique is that there can be considerable damage to the follicular unit as they are cored out from the scalp, one at a time because the skills to accomplish this are difficult to learn and master and the use of a sharp or dull punch to accomplish this runs into problems produced by the body's collagen which take on different consistencies with different people."}
-{"text": "The present invention relates to a magnet device for animals, in particular cattle, and particularly to a magnet device adapted to be arranged in the stomach of cattle and to collect scrap iron pieces taken down in the stomach so that the scrap iron pieces are prevented from entering into the intestines.\nCattle grazing in meadows sometimes swallow grass nails, wires or the like falling on the ground. These scrap iron pieces remaining in the stomach are slightly dangerous for cattle, but when they enter into the intestines from the stomach, it causes a great danger which can be fatal to the cattle.\nIn order to prevent such a danger, there has been proposed a device which consists of a magnet formed in a bar or cylinder at a size to not enter into the intestines from the stomach and which is adapted to remain in the stomach of cattle thereby collecting scrap iron or the like invading into the stomach.\nHowever, since such a device of a bar or cylinder magnet has magnetic pole portions limited to both end portions thereof, portion for performing a magnetizing function which collects the scrap irons is also limited to only both the end portions and its collecting effect has been insufficient.\nIn short, since the magnet device of this kind must remain in the stomach of cattle after the cattle has once swallowed it, without removing it out of the body of the cattle during its life, it was necessary to have the wide portion to magnetize scrap irons and to have a strong magnetizing force to magnetize the scrap irons.\nThe present invention is developed in the light of the above circumstances and in order to eliminate the above defects of the prior device, and contemplates to provide a magnet device which has the magnetizing portion not only at both ends but at peripheral wall portion and which has a magnetizing force intensified by concentrating flux of magnetic pole.\nAccording to the present invention, a magnetic device is constituted by a plurality of magnets in a short cylinder shape having magnetic poles of polarity differing from each other at both end surfaces and piled in a longitudinal direction through the intermediary of at least a magnetic plate between each magnet, magnetic poles of the magnet opposing through the intermediary of the magnetic plate being the same polarity whereby the flux concentrates at the magnetic plate (yoke) portion and the magnetizing force increases.\nFurther, according to the invention, a magnet device is constituted by a plurality of magnets in a short cylinder shape having those outer diameter which are similar to or smaller than the outer diameter of magnetic plate, in order to decrease leakage flux and to effectively utilize the flux.\nOther objects and features of the invention will appear in the following description taken with reference to the accompanying drawing."}
-{"text": "1. Field of the Invention\nThe present invention relates to a plasma welding method for a golf club head. Particularly, the present invention relates to the plasma welding method of tilting a plasma nozzle a predetermined angle with respect to the golf club head for avoiding obstruction of a neck portion or a hosel of the golf club head. More particularly, the present invention relates to the plasma welding method for welding on a connection line, within a predetermined section, between a striking plate and a club head body.\n2. Description of the Related Art\nReferring now to FIGS. 1 and 2, a conventional golf club head includes a club head body 91 and a striking plate 92 mounted to the club head body 91, as disclosed in TWN Patent Pub. No. 585,792, entitled \u201cgolf club head and manufacturing method therefor (II),\u201d TWN Patent Pub. No. I225,421, entitled \u201cgolf club head structure,\u201d TWN Patent Pub. No. I226,251, entitled \u201cconnecting structure for a striking plate with a golf club head,\u201d and U.S. Pat. No. 6,099,414, entitled \u201cgolf club head and method for producing the same\u201d etc. A connection line 93 exists between the club head body 91 and the striking plate 92 which are welded by a suitable welding method. In welding operation, a variety of welding methods, such as tungsten inert gas (TIG) welding, laser welding, plasma welding or other suitable welding methods, may be selectively used.\nThe above plasma welding may be selected to weld the club head body 91 and the striking plate 92, which carries out advantages of high-energy, high-accuracy, high-speed and high-quality welding, and further minimizes or attenuates a heat affected zone of material in a welding process.\nIn plasma welding operation, a movable plasma nozzle 94 is disposed on a slide track (not shown) and aligned with the connection line 93 existing between the club head body 91 and the striking plate 92. Subsequently, the plasma nozzle 94 is moved downward to an operating level above the connection line 93, and welded along the connection line 93 by an automatic welding process. The club head body 91 has a neck portion 91 and a hosel 912 connected thereto. It should be noted that the neck portion 911, however, is located between the hosel 912 and the club head body 91 and extended upward from a level of the striking plate 92. The neck portion 911 may unavoidably obstruct or interfere in a runway of the plasma nozzle 94 in the plasma welding since the neck portion 911, to a certain extent, is too close to the connection line 93. Consequently, the plasma nozzle 94 running at the operating level cannot pass through above a section of the connection line 93 adjacent to the neck portion 911 of the club head body 91 due to interference with the neck portion 911.\nAs explained above, the automatic plasma welding process is only suitable for welding the other section of the connection line 93 away from the neck portion 911 of the club head body 91. Meanwhile, other suitable welding methods, a TIG welding method for example, may be used to manually weld the section of the connection line 93 adjacent to the neck portion 911 of the club head body 91. Such welding method may cause poor quality, increased deformation and low efficiency of welding on the connection line 93 adjacent to the neck portion 911 of the club head body 91. Accordingly, such practice, however, may limit the welding process to carry out to be a complete automatic welding process. Hence, there is a need for improving the plasma welding process for the golf club head.\nThe present invention intends to provide the plasma welding method of tilting a plasma nozzle a predetermined angle with respect to a golf club head for avoiding obstruction of a neck portion or a hosel of the golf club head in such a way to mitigate and overcome the above problem. This permits the plasma nozzle passing through a connection line adjacent to the neck portion of the golf club head in the plasma welding process. Accordingly, the plasma welding method carries out a complete automatic plasma welding process so as to improve the welding efficiency and quality."}
-{"text": "Poly(vinyl chloride) (\u201cPVC\u201d) is commonly utilized as a material of construction for consumer and industrial goods. PVC possesses the advantageous features of low cost, durability, moisture resistance, tailored stiffness, dimensional stability, and flame retardancy. Virgin PVC is readily processed into sheet, tubes, and other forms using conventional processing equipment such as extruders and thermal compression bonding equipment. The reported density of poly(vinyl chloride) is 1.45 g/cm3.\nPlasticizers can be added to PVC, which is known to be rigid in the absence of such additives, to make it more flexible and more suitable for an even broader range of applications, including but not limited to applications such as plumbing and electrical cable insulation. The addition of plasticizers and other such flexibilizing agents will lower the modulus and the density of the PVC.\nThe use of chemical and physical inert gas blowing agents can also be used to lower the density of PVC substrates.\nWhile PVC is attractive for many first life commercial uses, the opportunities to reclaim and reuse PVC for subsequent future life applications are limited. The recycling of reclaimed PVC possesses inherent challenges due both to the processing difficulties caused by the various and often unknown additives and fillers that may have been used during the first life and to the sorting and separation of reclaimed PVC to generate a more homogenous raw material for use in future applications. Specifically, a significant amount of reclaimed PVC sources incorporate mixtures of other materials or components, for example pigments, colorants, fillers, plasticizers and the like, that limit PVC's potential reuse in future applications. Reclaimed PVC materials introduce feedstock variability for which conventional melt processing may be ill-suited and which limits the commercial utility for future life applications of PVC. This challenge becomes more acute as one tries to use a high fraction of reclaimed PVC within conventional processing methods."}
-{"text": "1. Field\nThe present disclosure relates to computer systems and methods in which data resources are shared among data consumers while preserving data integrity and consistency relative to each consumer. More particularly, the disclosure concerns implementations of mutual exclusion mechanisms such as reader-writer locking.\n2. Description of the Prior Art\nBy way of background, reader-writer synchronization is a mutual exclusion technique that is suitable for use in shared memory multiprocessor computing environments to protect a set of shared data. One type of reader-writer synchronization, known as reader-writer locking, allows read operations (readers) to share lock access in order to facilitate parallel data reads, but requires write operations (writers) to obtain exclusive lock access for writing the data. The technique is well suited to shared memory multiprocessor computing environments in which the number of readers accessing a shared data set is large in comparison to the number of writers, and wherein the overhead cost of requiring serialized lock acquisition for readers would be high. For example, a network routing table that is updated at most once every few minutes but searched many thousands of times per second is a case where serialized read-side locking would be quite burdensome.\nReader-writer locks are conventionally implemented using a single global lock that is shared among processors. This approach requires readers and writers to contend for one global lock on an equal footing, but produces memory contention delays due to cache line bouncing of the lock between each processor's cache. Insofar as reader-writer locks are premised on the existence of a read-intensive processing environment, readers may be unduly penalized, especially if their critical sections are short and their lock acquisition frequency is high. A distributed reader-writer lock approach is presented in Hsieh and Weihl, \u201cScalable Reader/Writer Locks for Parallel Systems\u201d, 1991. It requires the readers to acquire only a local per-processor reader/writer lock that will usually reside in the memory cache of the processor that hosts the acquiring reader. However, the writers must acquire all of the local reader/writer locks, which degrades writer performance due to memory contention, and in some cases due to new readers being allowed to starve a writer while the latter is waiting for one of the local reader/writer locks. A further disadvantage associated with both non-distributed and distributed reader-writer locking is that lock acquisition imposes a burden on readers, even in the absence of a writer. Reader-writer locks are typically implemented as semaphores, mutex locks and spinlocks. Acquiring each of these lock types often imposes the cost of atomic instructions and/or memory barriers. In a read-mostly computing environment, the overhead associated with these operations falls mostly on readers.\nImproved read-side performance is provided by the locking technique disclosed in commonly-owned U.S. Pat. No. 7,934,062, which requires no read-side lock acquisition except when a writer announces its intention to acquire the reader-writer lock. However, the write-side performance of this method can be degraded in systems with many processors. This is because writers must wait for a grace period to elapse before acquiring the reader-writer lock. All processors must pass through a quiescent state that guarantees each reader will have an opportunity to note the writer's locking attempt, and thereby synchronize on the reader-writer lock.\nThe present disclosure introduces techniques for reducing writer latency in large multiprocessor systems that employ data synchronization mechanisms, such as the grace period-based reader-writer locking approach disclosed in U.S. Pat. No. 7,934,062 or the distributed locking scheme proposed by Hsieh and Weihl. A technique for reducing writer latency in a multithreaded user-mode embodiment of the Hsieh and Weihl distributed locking method is also disclosed. The techniques disclosed herein are also useful for other synchronization operations, such as expedited grace period detection in multiprocessor systems implementing read-copy update (RCU) synchronization."}
-{"text": "1. Field of the Invention\nThe present application relates to an electronic apparatus provided with a slot to which an expansion module can be attached detachably.\n2. Description of Related Art\nRecent notebook computers are provided with a random access memory (RAM) that temporarily stores programs or data in order to execute various information processes at a central processing unit (CPU). As the size of a RAM increases, the size of an executable program or the number of simultaneously executable programs can be increased.\nTypically, a RAM is formed by mounting a memory chip on a board in the form of a module. When this memory module is attached to a memory slot provided in a notebook computer, the memory size of the notebook computer can be increased. This memory slot typically is sealed with a cover in order to prevent a foreign particle from entering from the outside. The cover is typically in the shape of a plate. Patent Documents 1 and 2 disclose plate-shaped covers. The covers disclosed in Patent Documents 1 (JP 2003-076439A) and 2 (JP H10-268976A) are attached to an apparatus main body through claw coupling or screw threading.\nHowever, in the case of a plate-shaped cover, when opening a memory slot in order to attach or detach a memory module, it is necessary to disconnect the claw coupling or screw threading, and then cause the cover to be detached by its own weight from an apparatus main body. More specifically, the electronic apparatus is initially positioned such that a face having the memory slot is oriented upward, claw coupling or screw threading is then disconnected, the electronic apparatus is then positioned such that the face having the memory slot is oriented downward, and the cover is detached by its own weight from the electronic apparatus. Subsequently, the memory module is attached or detached in a state where the face having the memory slot is oriented upward. Accordingly, when detaching a cover in order to attach or detach a memory module, it is necessary to change the posture of the electronic apparatus a plurality of times, which is very troublesome."}
-{"text": "1. Field of the Invention\nThe present invention relates to an ESD (electrostatic discharge) protection component and a method for manufacturing the EST protection component.\n2. Related Background Art\nThere is a known ESD protection component provided with an element body in which a plurality of insulator layers are stacked, a coil constructed by connecting a plurality of internal conductors to each other and arranged in the element body, and an ESD suppressor arranged in the element body and configured including first and second discharge electrodes arranged as separated from each other (e.g., cf. Japanese Patent Application Laid-Open Publication No. 2003-123936 (which will be referred to hereinafter as Patent Literature 1)). There is another known ESD protection component provided with an ESD suppressor configured including first and second discharge electrodes arranged as separated from each other, and a discharge inducing portion kept in contact with the first and second discharge electrodes so as to connect mutually opposed portions of the first and second discharge electrodes to each other and containing metal particles, in which a cavity portion is arranged so as to be in contact with the foregoing mutually opposed portions of the first and second discharge electrodes and with the discharge inducing portion (e.g., cf. Japanese Patent Application Laid-Open Publication No. 2011-243896 (which will be referred to hereinafter as Patent Literature 2))."}
-{"text": "The invention relates to communication between processors.\nMulti-processor computer systems have more than one processor. Each processor executes a separate stream (xe2x80x9cthreadxe2x80x9d) of instructions. It is sometimes necessary for two processors of a computer system to communicate data between themselves.\nIn one general aspect of the invention, a method of communicating between a first and a second processor includes the first processor sending a datum over a common control bus, and the second processor receiving the datum from the common control bus.\nAdvantages and other features of the invention will become apparent from the following description and from the claims."}
-{"text": "The present invention is directed to reflectors used in the radiant section of a fired heater, and more particularly to radiant reflectors provided on a refractory wall centered in the spacing between the radiant tubes.\nCombustion equipment is generally operated in chemical plants, petrochemical plants and refineries. The equipment may include industrial heaters, furnaces or plant boilers. This equipment is generally designed with bare or smooth-walled tubes, or with partially studded tubes as disclosed in my earlier U.S. Pat. No. 6,364,658, which is hereby incorporated herein by reference in its entirety. Use of tubes in radiant sections usually exposes the front half of the tube to direct flame radiation, while limiting the exposure of the rear half or dark side of the tube to reflected radiation.\nThe heat flux distribution around the circumference of a conventionally fired tube at a conventional spacing of two tube diameters is depicted in FIG. 1. A flame or radiating plane is on one side of the tube and a refractory wall is on the other. The front half of the tube surface faces the flame (point A) and receives a higher heat flux as compared to the rear half facing the refractory wall (point B). Point A receives heat flux only from direct flame radiation, while point B, facing the refractory wall, receives only reflected radiation coming from the refractory wall. Points between point A and point B receive varying amounts of both direct and reflected radiation, depending upon their location along the tube.\nThe standard distance between tubes is two tube diameters from center-to-center, and 1.5 diameters from the center of the tubes to the refractory wall, for most operations in the chemical and petrochemical industries, as shown in FIG. 2. The heat flux distribution in FIG. 1 is based on this configuration. For the purposes of an illustration using fluxes typical in a conventional fired heater, where the highest heat flux at point A is 18000 Btu/hr-ft2, the diametrically opposed counterpart (point B) receives only 6000 Btu/hr-ft2. The rear half of the tube transfers only 24% of the total heat absorbed by the tube; this includes both the direct and reflected radiation, as seen in FIG. 3. The average flux for the tube amounts to 10,000 Btu/hr-ft2.\nMore than 85% of the heaters in the industry have such a large flux differential between the front and the rear side of the tube, as this illustration depicts. A significant compromise is made on the overall heat-receiving capacity of the tube in order to keep the flame-front side (point A) within safe working temperatures.\nTo make the heat flux distribution in the tube more uniform, one approach of the furnace designers has been to increase the center-to-center tube spacing requirements from 2 to 3 tube diameters. This design increases the flux at point B of the tube from 6,000 Btu/hr-ft2 to 9,000 Btu/hr-ft2 as shown in FIGS. 4A and 4B. The increased spacing has the beneficial result of increasing the heat-receiving capacity of the rear half of the tube for the 3D-spaced tubes, while heat flux distribution on the front half of the tube is generally the same as for the 2D-spaced tubes. This results in an increase of the average heat flux to 12,000 Btu/hr-ft2 for the entire tube. However, the drawback of this solution is apparent. With an increase in tube spacing there is a corresponding increase in the size of the heater. This increases the cost and space requirements for the heater.\nAnother prior art approach improves the heat flux distribution by placing radiating flames on opposing sides of the tubes in a so-called xe2x80x9cdouble-firedxe2x80x9d design. A comparison is shown between one radiating flame (A) and two radiating flames (B) in FIGS. 5A and 5B, respectively. This design is commonly used in chemical processes that mandate a more uniform heat flux distribution, such as, for example, in delayed cokers, high-pressure hydrotreaters, ethylene furnaces, and the like. In a double-fired system, the front (point A) and rear (point B) portions of the tube have the same heat flux rate due to direct flame radiation, and the points at the margins between the front and rear receive relatively less direct flame radiation. The corresponding distribution of the heat flux, for the illustrative example, is 18,000 Btu/hr-ft2 for the front and the rear locations, 13,500 Btu/hr-ft2 at the margins between the front and rear faces, i.e. the middle area of the tube (point M at the 90 and 270 degree positions), resulting in an average flux of 15,000 Btu/hr-ft2. The double-fired design brings with it the disadvantage that the heater has to be much larger, as much as twice the size as a single-fired unit, and correspondingly more expensive.\nThe present state of technology for heaters with a standard spacing of 2 tube-diameters will have a relative flux ratio of 1 to 1.8 between the average flux and the maximum flux, whereas a heater with a 3 tube-diameter spacing will have a relative flux ratio of 1 to 1.5, as shown in API Standard 530, Calculation of Heater-Tube Thickness in Petroleum Refineries, American Petroleum Institute (1988), Figure C-1 Ratio of Maximum Local to Average Heat Flux Curves, page 103.\nThe 3 tube-diameter design is less common in the industry and the vessel must be significantly larger than a 2 tube-diameter design. The average to maximum flux ratio of the double-fired tubes is significantly lower at 1 to 1.2, but is a more costly alternative of the three designs for an industrial plant.\nA recent improvement in the flux distribution as described in my U.S. Pat. No. 6,364,658, involves the placement of extended surfaces such as studs or fins on the dark side of the tubes in a single-fired arrangement. This improves the heat transfer to the dark side of the tubes primarily by increasing the convection heat transfer. Still, in the standard tube arrangement with smooth walls, it is well known that 65.8% of the radiant heat from the flame is absorbed by the tubes, primarily the front half of the tubes facing the flame, and 34.2% goes through the spaces between the tubes to the refractory wall. The same percentages apply to the reflected radiation from the refractory onto the dark side of the tubes, i.e. 65.8% of the 34.2% is re-radiated to the rear half of the tubes, or 22.5%. In other words, 88.3% is absorbed by the tubes, front and back, and the balance of 11.7% is radiated back to the flame through the spaces between the tubes. It would be very desirable if a significant portion of this 11.7% could be directed onto the tubes instead of the flame. There thus remains a need for making the flux distribution even more uniform and/or for increasing the rate of heat absorption by the tubes.\nThe present invention utilizes radiation reflectors positioned on the refractory wall of a furnace, preferably in the spaces between the radiant tubes. The radiation reflectors provide surfaces which are angled, with respect to generally flat or curvilinear refractory surfaces behind the tubes, to reduce the radiation that is reflected between the tubes and increase the radiation reflected onto the dark side of the tubes. The use of the radiation reflectors thus increases the radiant flux delivered to the dark side of the tubes, increasing heat absorption and decreasing the ratio of the maximum to average flux. The radiation reflectors can also enhance convection heat transfer to the dark side of the tubes by increasing the velocity of the flue gases between the tubes and the refractory wall, thereby increasing the convection heat transfer.\nIn one aspect, the present invention provides radiation reflectors for use in a fired furnace comprising a plurality of parallel tubes arranged in a row between a flame on a radiant side and a generally flat or curvilinear refractory surface on a dark side. The radiation reflectors have a longitudinal base for abutment against the refractory surface. The base has opposite edges at either side thereof. A longitudinal cusp is opposite the base, and longitudinal reflective surfaces extend from each edge of the base to the cusp. The reflective surfaces have concavity in a plane transverse to a longitudinal axis, preferably parabolic sections in the transverse plane. An anchoring pin can extend transversely through each radiation reflector from the cusp into a subjacent structure.\nIn another aspect, the invention provides a fired furnace for heating petroleum, petrochemicals or chemicals. The furnace has a plurality of parallel tubes each disposed in a row between a flame on a radiant side thereof and a refractory surface on a dark side thereof. There are spaces between adjacent tubes. Radiation reflectors are positioned on the refractory surface opposite the spaces to reflect incident radiation from the flame away from the spaces and onto the dark side of the tubes. A central longitudinal bore is provided through each tube for the passage therethrough of a fluid to be heated. The row of tubes can be straight or circular. The radiation reflectors can be disposed longitudinally on either side of a flat surface of the refractory surface opposite a tube.\nIn a further aspect, the invention provides an improvement in a fired furnace. The furnace includes a plurality of parallel tubes disposed between a flame and a refractory wall. Adjacent tubes define a space between the tubes, and each tube includes a central longitudinal bore for the passage therethrough of a fluid to be heated and an outside diameter having a radiant side for exposure to radiation from the flame and a dark side essentially free of direct exposure to the flame. The improvement comprises positioning the radiation reflectors described above on the refractory wall opposite each space. Preferably, the reflective surfaces are parabolic sections in the transverse plane focused on the dark side of the adjacent tubes.\nIn a still further aspect of the invention, there is provided a method for improving the heat transfer in a fired furnace comprising a plurality of parallel tubes disposed between a flame and a refractory wall. Adjacent tubes define spaces between the tubes. The refractory wall comprises a generally flat or curvilinear surface opposite the tubes and spaces. The method includes the step of installing the radiation reflectors described above on the refractory wall opposite the spaces. The installation can include pinning the radiation reflectors with a pin extending from the cusp into the refractory wall. The radiation reflectors are preferably focused to reflect incident radiation from the flame onto the adjacent tubes on either side of a respective space. The tubes can have extended surfaces at least on the dark side. Where the tubes have smooth outside walls, the method can also include removing the smooth-walled tubes from the furnace and replacing them with tubes that have extended surfaces on a dark side opposite the refractory."}
-{"text": "Radio frequency RF) receivers are used in a wide variety of applications such as television receivers, cellular telephones, pagers, global positioning system (UPS) receivers, cable moderns, cordless phones, satellite radio receivers, and the like. As used herein, a \u201cradio frequency\u201d signal means an electrical signal conveying useful information and having a frequency from about 3 kilohertz (kHz) to hundreds of gigahertz (GHz), regardless of the medium through which such signal is conveyed. Thus an RF signal may be transmitted through air, free space, coaxial cable, fiber optic cable, etc. One common type of RF receiver is the so-called superheterodyne receiver. A superheterodyne receiver mixes the desired data-carrying signal with the output of tunable oscillator to produce an output at a fixed intermediate frequency (IF). The fixed IF signal can then be conveniently filtered and converted down to baseband for further processing. Thus a superheterodyne receiver requires two mixing steps.\nModern integrated circuit technology has allowed many of the circuits used in RF receivers to be combined on-chip and thus to substantially reduce the cost of the RF receiver. However this level of integration creates other problems. For example, signals from one part of the chip may be electrically or magnetically coupled to circuits in another part of the chip. These unwanted signal couplings can distort the desired signal and create artifacts that can be perceived by the viewer or listener. Traditionally, integrated circuit designers have used layout strategies to reduce coupling between circuits, such as physical separation, the addition of ground rings, a reduction in the length of conductors, etc. However these techniques, while still useful, are unable to completely eliminate the deleterious effects of electrically or magnetically coupled energy.\nThe use of the same reference symbols in different drawings indicates similar or identical items."}
-{"text": "1. Field of the Invention\nThe present invention relates to composite materials technology, and more specifically to a relatively light-weight, inexpensive, durable, high performance structural laminate composite material for use to 1000xc2x0 F., and above, which can advantageously be used in high temperature environments. More particularly, the preferred embodiment of the present invention relates to a graphite-fiber/phenolic-resin composite material which retains relatively high strength and modulus of elasticity at temperatures as high as 1,000xc2x0 F. (538xc2x0 C.). The material costs only 5 to 20 percent as much as refractory materials do. The fabrication of the composite includes a curing process in which the application of full autoclave pressure is delayed until after the phenolic resin gels. This modified curing process allows moisture to escape, so that when the composite is subsequently heated in service, there will be much less expansion of absorbed moisture and thus much less of a tendency toward delamination. In contrast, internal pressure caused by the expansion of moisture absorbed in other prior art composite materials like prior art graphite/epoxies and prior art graphite/polyimides causes delamination at temperatures in the range of 500 to 700xc2x0 F. (260 to 370xc2x0 C.).\n2. General Background\nAt the request of NASA/MSFC, Martin Marietta Manned Space Systems has performed an extensive development/verification activity for a composite nose cone for the external tank (ET). At the time of the initiation of this effort, there was no materials technology available to provide a nose cone which could withstand the high heating and structural loading of the ET nose cone without (a) requiring the use of secondary heat shield materials, (b) increasing the weight of the existing nose cone, and (c) significantly increasing the cost over the existing nose cone cost. There were high temperature polymeric composite materials available; however, none met all requirements. Carbon/phenolic laminates have been proven in rocket nozzle applications to be able to withstand extreme heating conditions; however, these materials did not possess the specific strength and stiffness required for a weight-effective structure. Also, recent data shows that the materials on the market today have the potential to xe2x80x9cply lift,xe2x80x9d or delaminate due to internal pressure caused by absorbed moisture, at about 500xc2x0 F. Graphite/polyimide laminates showed promising mechanical properties, but suffered from the moisture-induced delamination problem (also known as xe2x80x9cthermal shockxe2x80x9d) at temperatures below 700xc2x0 F. in laminates of the thickness required for a composite nose cone. Other technologies such as ceramic matrix composites and carbon/carbon were considered too expensive for this application. Therefore, a program was initiated to develop laminate material which could meet all requirements.\nU.S. Pat. No. 3,724,386 for xe2x80x9cAblative Nose Tips and Method for their Manufacturexe2x80x9d discloses in Example II heating graphite yarn impregnated with phenolic resin slowly to 160xc2x0 F. to slowly evaporate solvent from the resin (see column 8, lines 16-18).\nU.S. Pat. Nos. 4,100,322 and 4,215,161 for xe2x80x9cFiber-Resin-Carbon Composites and Method of Fabricationxe2x80x9d disclose impregnating graphite yarn with phenolic resin under vacuum and a temperature of about 150xc2x0 F. until the solvent has gone and the resin gels, then further heating the composite to cure it. However, the solvent stripping process was interrupted twice and each time pressure of 200 psig was applied to the composite material. It is then subjected to pyrolysis, and then pores of the composite are impregnated with phenolic resin. After this, the phenolic resin is cured at about 350xc2x0 F. The resulting structure is said to be graphite/carbon/phenolic composite, and its porosity is disclosed to be 4%. A carbon/carbon/phenolic composite described therein is said to have a porosity of 5.8%.\nU.S. Pat. No. 4,659,624 for xe2x80x9cHybrid and Unidirectional Carbon-Carbon Fiber Reinforced Laminate Compositesxe2x80x9d discloses a method similar to the method disclosed in U.S. Pat. Nos. 4,100,322 and 4,215,161 (and with similar materials), but one in which more resin is added and pyrolized up to 5 times. This patent points out at column 2, line 50 through column 3, line 2 that it is important to properly initially cure laminate materials to provide interconnecting pores which allow the escape of gases formed during post-cure pyrolysis.\nU.S. Pat. No. 4,957,801 for xe2x80x9cAdvance Composites with Thermoplastic Particles at the Interface Between Layersxe2x80x9d discloses a resin-impregnated fiber layer with outer layers of resin thereon. The fiber can comprise, for example, graphite.\nU.S. Pat. No. 5,288,547 for xe2x80x9cToughened Resins and Compositesxe2x80x9d discloses a composite in which a porous membrane film of thermoplastic material is sandwiched between two layers of resin-impregnated fibers, and then the composite is cured in an autoclave, for example. The resin can be, for example, phenolic resin.\nU.S. Pat. No. 5,359,850 for xe2x80x9cSelf Venting Carbon or Graphite Phenolic Ablativesxe2x80x9d discloses a resin-impregnated reinforcing cloth made of, for example, graphite fibers with degradable fibers interwoven therewith. The degradable fibers are chosen such that they degrade at a temperature of about 400xc2x0 F. to 500xc2x0 F. so that they will provide passageways for the gaseous decomposition products produced as the resin matrix approaches the char temperature. In this patent, foreign material is introduced to create porosity. The fabric weave is altered by introducing a low-temperature degradable thread which may not assure fabric strength properties. The porosity which is created by this process is uniform. There is a definite pattern when the foreign material is replaced by voids. It is believed that the addition of these special degradable fibers will add to the cost of the material. Further, it is believed that in some cases the degradable fibers might not burn away before the plies blow apart.\nU.S. Pat. No. 5,360,500 for xe2x80x9cMethod of Producing Light-Weight High-Strength Stiff Panesxe2x80x9d discloses a panel made by a pair of surface members separated and supported by an internal core in which spaces or interconnected pores provide vents to an edge of the panel so that gas can flow through the vents during a pyrolysis process. The vents are on the order of 10 mm in diameter.\nNone of these patents discloses a composite material with a weight, thickness, structural performance, and pore structure as advantageous for use in a nose cone of the external tank of the space shuttle, or other high temperature structural applications, as the material of the present invention.\nA novel materials technology has been developed and demonstrated for providing a high modulus composite material for use to 1000xc2x0 F. The material of the present invention can be produced at 5-20% of the cost of refractory materials, and has higher structural properties. This technology successfully resolves the problem of xe2x80x9cthermal shockxe2x80x9d or xe2x80x9cply lift,xe2x80x9d which limits traditional high temperature laminates (such as graphite/polyimide and graphite/phenolic) to temperatures of 550-650xc2x0 F. in thicker (0.25xe2x80x3 and above) laminates. The technology disclosed herein is an enabling technology for the nose for the External Tank (ET) of the Space Shuttle, and has been shown to be capable of withstanding the severe environments encountered by the nose cone through wind tunnel testing, high temperature subcomponent testing, and full scale structural, dynamic, acoustic, and damage tolerance testing.\nIn the present invention, cure conditions (temperature, pressure, vacuum) and cure apparatus (specific vacuum bag methodology) are manipulated to produce a graphite/phenolic composite laminate with a permeable microstructure comprising an interconnected network of pores which allows moisture to escape from the composite material when the composite material is heated; this helps prevent delamination (xe2x80x9cply liftxe2x80x9d or xe2x80x9cthermal shockxe2x80x9d) when the material is heated to temperatures above 500xc2x0 F. The graphite/phenolic composite of the present invention can be used for components for applications requiring high strength and stiffness upon exposure to very high heating (e.g. rocket nozzles for missiles or launch boosters, fire walls, heat shields, circuit boards, secondary structure on missiles or launch vehicles which see high aerodynamic heating, and parts to be used on the leading edge of aerodynamic products (airplanes, jets, rockets, fuel tanks for aerospace structures, etc.)).\nThe present invention comprises a method of producing a composite material, comprising the steps of:\nimpregnating a fiber material with a resin to create a resin-impregnated fiber material;\nwithout applying pressure, heating the resin-impregnated fiber material under vacuum at a sufficient temperature for a sufficient amount of time until the resin gels; and\napplying temperature (and, optionally, pressure) for a sufficient period to cure the resin-impregnated fiber material. The starting percentage by weight of fiber material (before being cured) is preferably 30-80%, with the balance resin. The resulting porosity of the composite material is preferably at least 3% by volume, more preferably about 3-25% by volume, and most preferably about 7-14% by volume.\nThe preferred embodiment of the method of the present invention of producing a composite material comprises the steps of:\n(i) impregnating a graphite fiber material with a phenolic resin to create a resin-impregnated fiber material, in a ratio of 30-80% by weight graphite fiber and 20-70% by weight phenolic resin;\n(ii) placing the resin-impregnated fiber in an autoclave or oven;\n(iii) applying full vacuum and/or pressure;\n(iv) raising the temperature to cause the resin to flow and initiate cure,\n(v) holding the material at a temperature to allow gellation of resin while volatiles are being released;\n(vi) raising the temperature for final cure if required;\n(vii) cooling the material;\n(viii) removing the material from the autoclave or oven;\n(ix) post-curing the composite laminate material removed from the autoclave, if required.\nThe present invention includes the composite material made by the method of the present invention disclosed herein, as well as a composite material, produced by any method, having a composition and structure which is the same as the composite material produced by the method of the present invention disclosed herein.\nThe material of the present invention comprises a high performance structural laminate composite material for use in high temperature applications, consisting essentially of resin-impregnated fiber, the resin-impregnated fiber consisting essentially of:\n(a) preferably 50-80% by weight fiber, and\n(b) preferably 20-50% by weight cured resin, the composite material having:\n(c) a permeability sufficient to allow moisture to escape from the composite material, without causing plylift, when the composite material is heated to temperatures up to 1000xc2x0 F. More preferably, the permeability is sufficient to allow moisture to escape from the composite material, without causing plylift, even when the composite material is heated to temperatures above 1000xc2x0 F. The material of the present invention has a microscopic construction which provides permeability that is sufficient to allow moisture to escape therefrom as it is heated to temperatures up to 1000xc2x0 F. and above without exhibiting ply-lift.\nThe composite material preferably has an across-ply permeability having a Darcys constant of at least 10xe2x88x9215 cm2. More preferably, the across-ply permeability of the composite material has a Darcy\"\"s constant of at least 10xe2x88x9214 cm2. Most preferably, the across-ply permeability of the composite material has a Darcy\"\"s constant of at least 10xe2x88x9213 cm2.\nThe material of the present invention comprises a high performance structural laminate composite material for use in high temperature applications, consisting essentially of phenolic resin-impregnated graphite fiber, the phenolic resin-impregnated graphite fiber consisting essentially of:\n(a) preferably 50-80% by weight graphite fiber; and\n(b) preferably 20-50% by weight cured phenolic resin, the composite material having:\n(c) a permeability sufficient to provide a network of pores which allows moisture to escape from the composite material, without causing plylift, when the composite material is heated.\nThe percentage by weight of graphite fiber is more preferably 60-80%, and the percentage by weight of cured phenolic resin is more preferably 20-40%. Most preferably, the percentage by weight of graphite fiber cloth is 65-75%, and the percentage by weight of cured phenolic resin is 25-35%.\nPreferably, the porosity is 3-25% by volume. Most preferably, the porosity is 7-14% by volume.\nPreferably, the compressive strength of the material after exposure to temperatures above 700xc2x0 F. for several minutes is at least 50% of the compressive strength of the material immediately after being cured, the shear strength of the material after exposure to temperatures above 700xc2x0 F. for several minutes is at least 50% of the shear strength of the material immediately after being cured, and the compressive strength of the material at 900xc2x0 F. is at least 25% of the compressive strength of the material at room temperature.\nThe graphite fiber cloth was selected to have a combination of high strength, high modulus, good thermo-oxidative stability, and moderate cost. The optimum fiber type to provide this balance is a fiber made from a polyacrylynitrile (PAN) precursor, such as the Toho G30-5001 fiber used in the development documented herein. Similar fibers are Hercules AS4 and IM-7, and Amoco T300, T650-35, and T650-45. Fiber types which were not selected were fibers based on pitch precursors (e.g., Amoco P-75 and P-100), or fibers based on rayon precursors. Pitch based fibers are much more expensive and do not have adequate strength. Rayon based fibers do not have the desired strength or modulus. The selected fiber was woven into an 8-harness satin fabric to facilitate part fabrication. The selected fiber can be, for example, an eight-harness fabric woven from Toho G-30/500-3K graphite fiber. The resin is advantageously selected from a group consisting of phenolics, bismaleimides (BMIs), polyimides, cyanate esters, epoxies, or any blend of these resins. The resin can comprise Cytec 506 phenolic resin.\nPreferably, the graphite fiber material has a minimum tensile strength of at least 300 KSI, more preferably at least 400 KSI, and most preferably at least 500 KSI, a minimum modulus of at least 20 MSI, more preferably at least 25 MSI, and most preferably at least 30 MSI, and relatively low cost. Most preferably, the resin is phenolic resin.\nThe material can consist of phenolic resin-impregnated graphite fiber cloth, and the phenolic resin-impregnated graphite fiber cloth can consist of graphite fiber cloth and cured phenolic resin.\nThe present invention also includes apparatus comprising a component which requires high strength and stiffness upon short term exposure to very high heating, made of the material of the present invention. The component can be a rocket nozzle, a part for an aerodynamic vehicle, or some other component exposed to high heating. The component can be part of a fire wall or heat shield.\nFurther, the present invention comprises vacuum bag apparatus for producing a composite laminate material having a network of pores. This vacuum bag apparatus can comprise:\n(a) a base for receiving the laminate material thereon;\n(b) a non-stick layer to be received on the laminate material for helping to prevent the laminate material from sticking to layers above the non-stick layer;\n(c) a first volatiles flow and resin retaining layer above the non-stick layer for allowing volatiles, but not the majority of the resin, to escape from the laminate material through the first volatiles flow and resin retaining layer as heat is applied and the vacuum is drawn in the bag apparatus;\n(d) a bleeder layer on the first volatiles flow and resin retaining layer for absorbing most of the resin which flows through the first volatiles flow and resin retaining layer;\n(e) a second volatiles flow and resin retaining layer on the bleeder layer for allowing volatiles, but very little resin, to flow through the bleeder layer as heat is applied and the vacuum is drawn in the vacuum bag apparatus;\n(f) a first gas-flow layer on the second volatiles flow and resin retaining layer for allowing gas to flow evenly through the vacuum bag apparatus when a vacuum is drawn in the apparatus;\n(g) a lateral gas-flow layer surrounding the laminate material to ensure that volatiles can flow out of the laminate in virtually any direction;\n(h) a vacuum bag layer attached to the base in an air-tight manner, the base and the vacuum bag layer enclosing the laminate material and the non-stick layer, the first volatiles flow and resin retaining layer, the bleeder layer, the second volatiles flow and resin retaining layer, the gas-flow layer, and the lateral gas-flow layer. In certain circumstances, one or more of the layers can be omitted, as described further below. The vacuum bag apparatus preferably also comprises a port in the vacuum bag layer communicating with a vacuum source for allowing a vacuum to be pulled in the bag. In a room at standard temperature and pressure, the vacuum causes a pressure of about 15 psi to be applied to the laminate in the vacuum bag. The application of additional pressure may not be a necessary step to make the present invention work.\nIt is an object of the present invention to provide a high-strength, low weight, high temperature material which has sufficient permeability to allow moisture to exit therefrom, even when the material has a thickness of more than 0.40xe2x80x3, when heated to temperatures of above 500xc2x0 F., without damaging the internal structure of the material.\nIt is an object of the present invention to provide a high-strength, low weight, high temperature material which has sufficient permeability to allow moisture to exit therefrom, even when the material has a thickness of more than 0.40xe2x80x3, when heated to temperatures of above 1000xc2x0 F., without damaging the internal structure of the material.\nIt is another object of the present invention to provide a method of making such material.\nA further object of the present invention is to provide components made of such material.\nIt is also an object of the present invention to provide a material which can withstand the high heating and structural loading of the ET nose cone without (a) requiring the use of secondary heat shield materials, (b) increasing the weight of the existing nose cone, and (c) significantly increasing the cost over the existing nose cone cost.\nAnother object of the present invention is to provide an ET nose cone made of this material.\nUnlike many prior art methods of producing composite material, in the method of the present invention, there is no pyrolizing step (the composite material of the present invention is not pyrolized). In the method of the present invention, unlike the method of U.S. Pat. No. 5,359,850: no foreign material is introduced to create porosity; the fabric weave is not altered and areal weight of fabric is constant, which assures strength properties; the porosity which is created by the process of the present invention is random and spread over the composite laminate; no material is decomposed by the method of the present invention; and the cost of creating the porosity is relatively low.\nBecause of the high permeability of the material of the present invention, it is believed by the inventors that there will be no ply lift at any thickness, whether the laminate is at least 0.1 inch thick, at least 0.2 inch thick, at least 0.4 inch thick, or even more than 4 inches thick.\nAlthough the specific examples described herein relate to graphite fiber and phenolic resin, other appropriate fibers and resins could be used in conjunction with the present invention."}
-{"text": "1. Field of the Disclosure\nThe subject disclosure relates generally to oilfield drilling, and more particularly to bottom hole assemblies and tools for orienting a bottom hole assembly (BHA).\n2. Background of the Related Art\nIn conventional drilling, the BHA is lowered into the wellbore using jointed drill pipes or coiled tubing. Often the BHA includes a mud motor, directional drilling and measuring equipment, measurements-while-drilling tools, logging-while-drilling tools and other specialized devices. A simple BHA having a drill bit, various crossovers, and drill collars is relatively inexpensive, costing a few hundred thousand US dollars, while a complex BHA costs ten times or more than that amount.\nMany drilling operations require directional control so as to position the well along a particular trajectory into a formation. Directional control, also referred to as \u201cdirectional drilling,\u201d is accomplished using special BHA configurations, instruments to measure the path of the wellbore in three-dimensional space, data links to communicate measurements taken downhole to the surface, mud motors, and special BHA components and drill bits. The directional driller can use drilling parameters such as weight-on-bit and rotary speed to deflect the bit away from the axis of the existing wellbore. In some cases, e.g. when drilling into steeply dipping formations or when experiencing an unpredictable deviation in conventional drilling operations, directional-drilling techniques may be employed to ensure that the hole is drilled vertically.\nDirection control is most commonly accomplished through the use of a bend near the bit in a downhole steerable mud motor. The bend points the bit in a direction different from the axis of the wellbore when the entire drill string is not rotating. By pumping mud through the mud motor, the bit rotates though the drill string itself does not, allowing the bit alone to drill in the direction to which it points. When a particular wellbore direction is achieved, the new direction may be maintained by then rotating the entire drill string, including the bent section, so that the drill bit does not drill in a direction away from the intended wellbore axis, but instead sweeps around, bringing its direction in line with the existing wellbore. As it is well known by those skilled in the art, a drill bit has a tendency to stray from its intended drilling direction, a phenomenon known as \u201cdrill bit walk\u201d. A device for addressing drill bit walk is shown in U.S. Pat. No. 7,610,970 to Sihler et al. issued Nov. 3, 2009, which is incorporated herein by reference.\nThe use of coiled tubing with downhole mud motors to turn the drill bit to deepen a wellbore is another form of drilling, one which proceeds quickly compared to using a jointed pipe drilling rig. By using coiled tubing, the connection time required with rotary drilling is eliminated. Coiled tube drilling is economical in several applications, such as drilling narrow wells, working in areas where a small rig footprint is essential, or when reentering wells for work-over operations.\nIn coiled tubing drilling, a BHA with a mud motor is attached to the end of a coiled tubing string. Typically, the mud motor has a fixed or adjustable bend housing in order to drill deviated holes. Because the coiled tubing is unable to rotate from surface, a so called orienter tool is used as part of the BHA to \u201corient\u201d the bend of the mud motor into the desired direction. There exists a multitude of different designs for the drive systems of such tools. Some designs support continuous rotation such as electric motor and gearbox drives, while others only permit rotation by a certain limited angle. The orienter tool is typically a high-torque, low-speed device, wherein the design of the drive system provides a torque output which can at least match the reactive torque exerted by the drilling mud motor.\nFor example, some orienter tools have utilized planetary gears in an effort to drive the output shaft. Basically, creating a torque on an output shaft means that a tangential force has to be exerted. By way of example, an output torque of 1,000 ft-lbs from a 2-inch diameter shaft means a tangential force of 12,000 lbs. This amount of force will quickly yield any material unless the tangential force is evenly distributed over a sufficient area to reduce the stress levels. In a conventional planetary stage with a size constraint on the order of 3 inches in diameter, the limits of how much bending force the gear teeth can take, and how much stress the planet carrier is capable of supporting will be much below 1000 ft-lbs of torque."}
-{"text": "There is a general need for a simple and accurate apparatus for rapidly distributing material in the form of small solid bodies simultaneously and uniformly into a plurality of spaced containers for further processing. For example, seeds sown in containers must be covered to prevent them from washing away, and to prevent the build-up of moss and algae on top of the containers. A number of devices have been developed to spread covering material, usually granite grit or basalt gravel, over the seed. These prior devices are relatively bulky, mechanically complex, and somewhat unreliable. For example, some of the prior devices employ revolving drum members, rotating discs, or shutter boxes, all of which are quite bulky and expensive, as well as not being satisfactorily precise. Consequently, there is a definite need for a simpler, less expensive, and more accurate manually operable spreader usable with various materials, which can be adapted for use with different types of containers and with different kinds of seed covering materials, or for distributing seeds themselves."}
-{"text": "The present invention relates generally to testing integrated circuits, and more particularly to built-in-test circuits and output response analyzers for integrated circuits.\nIntegrated circuits are typically tested multiple times while they are manufactured. Often, individual circuits are tested while they are part of a wafer, which contain thousands of integrated circuits. Nonfunctional die are identified, for example with an ink spot, during a test referred to as wafer sort. After wafer sort, the die are separated and packaged. The packaged devices are testing again\u2014this is referred to as final test. Additional testing may be done, for example sample devices may be tested under extreme environmental conditions.\nDuring these tests, test data, also referred to as test vectors, which typically include data and clock signals, are provided to the integrated circuit by a tester. The input test data may be generated by a circuit or software test pattern generator. Conventionally the integrated circuit operates on the input test data and provides output test data back to the tester. An output response analyzer in the tester checks the output test data for errors, and passes or rejects the device.\nIt is desirable to test each node in an integrated circuit. However, integrated circuits are becoming extremely complicated and may include hundreds of thousands of logic elements. At the same time, it is desirable to reduce the number of pins on the device in order to simplify device packaging and reduce printed circuit board complexity and space. The result is that many internal nodes on integrated circuits are difficult to reach electrically by device pins.\nAccordingly, it is desirable to include test circuitry on the integrated circuit itself, such that these internal nodes may be more thoroughly tested. Further, it is desirable to provide an internal test circuit that is capable of testing using test patterns other than simple all ones or all zeros patterns. Also, it is desirable to be able to perform such tests without the addition of complicated circuitry. It is also desirable that the internal circuitry require no or a limited number of pins, such that device pin count may be maintained."}
-{"text": "1. The Field of the Disclosure\nThe disclosure relates generally to a system and methods for determining whether an energetic substance or material has experienced a reaction (\u201cgo\u201d) or a non-reaction (\u201cno-go\u201d) for storage, transportation or in-process handling, and more particularly, but not necessarily entirely, to a system and methods using a video capturing device, a CPU or computer, sensitivity test equipment, such as an electrostatic discharge device or impact assessment device for testing and assessing the substance or material reaction or explosion sensitivities, and a set of rules or instructions to be followed for quantifying and determining whether a reaction has occurred or not.\n2. Description of Related Art\nReaction detection for sensitivity test equipment is not automated. As a consequence, the determination of whether a reaction actually occurred (i.e., a go reaction) or whether a no reaction occurred (i.e., a no-go reaction) is subjective in nature and is based on an operator's varying experience and perception of what actually transpired during a reaction event. Because of the subjectivity of determining whether a reaction occurred or not, standardization of sensitivity test results between laboratories is very difficult to achieve. Standardization between laboratories is possible provided an objective system and methods for determining whether a reaction occurred or not are used.\nIn the industry, various types of equipment are used to assist an operator in determining whether a reaction occurred or not while testing an energetic material or substance for sensitivity. Some of these operator assisted devices include, but are not limited to, a noise dosimeter, a gas analyzer, a light meter, a strain gauge, and a video capturing device (such as a standard or high-speed camera). One of the best ways in the industry to determine whether a reaction occurred or not is to use a high-speed video capture device for at least the following reasons: (1) it provides a reviewable visual record of the reaction event; and (2) it provides better spatial and temporal visual resolution of the reaction event.\nThe disclosure is directed to a unique and advantageous system and methods for determining whether a reaction occurred or not using automated equipment, thereby reducing the amount of subjectivity resulting from an operator determination based on an image or collection of images. The automated system and method may use a high-speed video capturing device, sensitivity test equipment, a computer processor and a set of rules or instructions that objectively compares a set of quantified image trial data to a set of quantified image baseline or background data to determine whether a reaction occurred.\nAn additional difference between the disclosure and what is done in the industry is instead of the operator making a subjective reaction determination based on sensory perception or a less subjective determination using a threshold for auditory, or a gas analyzer, or intensity of light, the operator makes a determination on quantifiable data relating to the acceptable level of error or the likelihood of a false positive and/or a false negative. With that level of error identified, the reaction detection threshold can then be calculated through the use of the system and method of the disclosure. Knowledge of the error or the likelihood of a false positive and/or a false negative allows for much better risk assessment of the sensitivity testing outcome.\nThe disclosure improves upon known techniques used in the industry by, inter alia, using four unique identifiers or quantifiers relating to the images identified as being of interest and significant. Those identifiers or quantifiers include: brightness, shape, buoyancy, and the uniformity of the event. The disclosure also detects decomposition or reaction of energetic materials. The method and process disclosed may also be advantageous in that multiple characteristics of the images are simultaneously quantified to determine if a reaction has occurred. It is noteworthy that none of the known systems or methods known to Applicant provides the above-identified advantages.\nThe features and advantages of the disclosure will be set forth in the description which follows, and in part will be apparent from the description, or may be learned by the practice of the disclosure without undue experimentation. The features and advantages of the disclosure may be realized and obtained by the methods and means of the instruments and combinations particularly pointed out in the appended claims. Any discussion of documents, acts, materials, devices, articles or the like which has been included in the present specification is not to be taken as an admission that any or all of these matters form part of the prior art base, or were common general knowledge in the field relevant to the disclosure as it existed before the priority date of each claim of this application."}
-{"text": "An image recognition and classification system (a machine vision system) includes a preprocessor in which a \"top-down\" method is used to extract features from an image, an associative learning neural network system which groups the features into patterns and classifies the patterns, and an attentional mechanism which focuses additional preprocessing and a neural network on relevant parts of an image.\nAttempts to recognize and classify images have led to construction of automated artificial machine vision systems and to development of strategies to learn patterns in images and to recognize and classify images by using the learned patterns. Those developing artificial systems have continually attempted to incorporate principles of biological systems into their strategies, because biological systems outperform all artificial systems, implemented or proposed, by a wide margin. For example, machine vision systems based on artificial neural networks have been implemented on digital parallel computers, but a parallel implementation only provides an increase in speed without an increase in performance. Thus, the goal of emulating the pattern recognition performance of biological systems still eludes computer scientists.\nIn order for a biological nervous system to discriminate objects two fundamental problems must be solved: object segmentation and binding. \"Object segmentation\" deals with distinguishing separate objects; \"binding\" deals with how specific attributes such as shape and depth, are linked to create an individual object. A question addressed by object segmentation mechanisms is to which overlapping object does a border belong? An image of an object may be occluded (divided) by an overlapping image, and will need to be reconstructed as a whole image. Models have been proposed to explain this process, for example, using artificial neural networks. (Sajda and Finkel, 1992)\nDuring the past half century, the theoretical infrastructure of machine vision systems has developed both top-down (beginning with large features of the image) and bottom-up (beginning at the lowest level of resolution, usually a pixel) views. However, actual development has focused almost exclusively on bottom-up approaches as exemplified by the title of Pentland's illuminating book From Pixels to Predicates, and comments therein such as: \"Processing is primarily data-driven (i.e., bottom-up), although it can be responsive to the goals and expectations at the higher levels.\" (Pentland, 1986, part 1, page 1).\nOngoing efforts have focused on the extraction of \"features\" in an image by local manipulations of small micro-features (often 3.times.3 rarely more than 9.times.9 pixel areas), with the intent of identifying larger features (macro-features) from their combination. The paucity of robust results from this approach may be attributed to several causes, two of the most important of which are (1) that the mathematical operations performed on the small areas are usually differential operators such as edge detectors that enhance rather than reduce noise; and (2) that not even humans are very good at visual recognition when allowed only a small instantaneous field of view. Similarity of tactual and visual picture recognition with limited field of view. Loomis et al. (1991).\nDuring this same time period, cognitive psychologists and neurobiologists have made impressive advances in research on the processing mechanisms that are at work in the visual cortex of mammals, particularly cats and monkeys. Electrophysiological and psychophysical experiments on cats and monkeys demonstrate a wide variety of feature selective cells in the visual cortex. In the mammalian cortex, these include simple cells (Hubel and Weisel, 1962), whose shape is closely approximated by a Gabor function (Daugman, 1985; Jones and Palmer, 1987) or a difference of Gaussian functions; end-stopped cells (often called first order hypercomplex cells) (Hubel & Weisel, 1965; Gilbert, 1977); color sensitive cells; and even cells that respond only to faces. (Desimone, 1991). Face-selective cells in the temporal cortex of monkeys. Desimone (1991).\nComplex cells and second order hypercomplex cells (Hubel and Weisel, 1962, 1965) are sensitive to the same features as simple and first order hypercomplex cells, respectively. One of the differences among these cells, of interest in the context of feature extraction from static images, is that the complex and second order hypercomplex cells have larger receptive fields than simple cells, and are insensitive to location of micro-features within their receptive fields.\nIn the development of artificial systems, preprocessing of data derived from an image has been used to extract features from an image and to select features for further processing by machine vision systems. Preprocessing generally proceeds in steps from the \"bottom-up,\" although \"top-down\" preprocessing has been suggested as a model for human vision. Preprocessing is accomplished by preprocessors, which may be implemented in hardware or software. In some systems, preprocessors have served as the first layer of a two layered neural network. Preprocessing strategies have included subdividing a whole image to be processed into sub-images. Various filters have been suggested to operate on the data, converting the data to a different form or value distribution. Control masks have been used to focus a network on a specific domain of an image.\nPrior approaches to the problem of modeling biological preprocessing have been addressed by Grossberg (1988) and Fukushima (1988). For example, the neocognitron neural network developed by Fukushima conceptually models simple, complex, first order and second order hypercomplex cells as well as layers of cells that are sensitive to higher order features. Second order hypercomplex cells are constructed from combinations of first order hypercomplex cells; complex cells are constructed from combinations of simple cells, and the like.\nOne means of making a complex cell insensitive to location, the approach used by Fukushima, is to design it to receive input from several adjacent simple cells, whose frequency and orientation tuning are similar. The complex cell is made sensitive enough to respond when only one of the simple cells responds to a stimulus. The result is a complex cell with the same frequency and orientation tuning as the simple cells, whose receptive field size is equivalent to the total receptive field size of all its input simple cells combined. Furthermore, the complex cell is insensitive to where in its receptive field the luminance pattern is located (i.e., the complex cell is insensitive to which simple cell has been activated.) Trying to apply biological principles to artificial vision systems, Porat and Zeevi (1989) determined from their work and the work of others, that \"primitives of image representations in vision have a wavelet form similar to Gabor elementary functions (EF's),\" and proposed a method for texture discrimination in images using a Gabor approach.\nAlthough Porat and Zeevi (1989) proposed that \"These localized operators (referring to Gabor functions) are also suitable for a pyramidal scheme of multiresolution which appears to be characteristic of vision, and can also serve as oriented-edge operators and in pattern recognition tasks,\" (p. 116), they adopted the prevailing approach to the process as a bottom-up hierarchy.\nAn alternative to extracting features using predefined, generally applicable fixed filters (detectors), such as generated by Gabor and end-stop filters, is to design a system that generates its own feature detectors. In biological systems, the feature detectors must be general enough to handle all possible inputs encountered during the life experiences of the animal. It has been shown that a linear neural network with a correlation rule, when stimulated by random noise, will develop feature detectors similar to the center-surround and Gabor filters found in some artificial visual systems. However, in most practical applications of artificial networks, the universe of possible inputs is more restricted. This suggests that a system for adaptive filter generation that can develop feature detectors specific to the range of images that are encountered in a practical application would be highly desirable. Self-modifying learning algorithms have been pursued wherein a learning algorithm learns about its own effectiveness and modifies itself so that it is the most effective algorithm for solving a certain class of problems.\nDespite extensive efforts and much progress, \"Forty years of research in artificial neural networks has yielded networks with the neural complexity of, perhaps, a sea slug.\" (Wenskay, 1991) Image recognition and classification remains a major frontier. The present invention advances toward this frontier."}
-{"text": "The invention relates to the field of computers, and more specifically to computer buses.\nPresent day computers comprise bus systems, onto which different devices may be plugged. More specifically, a bus system is often comprised of a bus controller and of a bus connected to the memory controller. Different devices may be connected to the bus, so as to be accessed by the bus controller.\nOne example of such buses is the DRAM bus designed by Rambus Inc. This bus is used for managing high speed DRAM devices. FIG. 1 is a schematic view of the architecture of this type of bus. It shows the memory controller 1, and the Rambus Channel 2. Several Direct Rambus DRAMs or Direct RDRams (trademark) 3-6 are connected to the Rambus Channel. As shown on FIG. 1, the memory controller as well as each of the RDRAMs comprises a Rambus interface 8 for using the bus. The bus 2 is terminated at one end by terminations, and is also connected to a reference voltage Vref as well as to a 400 MHz bus clock.\nAccording to the Rambus specification, there is also provided a power down mode; it is contemplated in the specification that the power down mode is used for reducing power consumption, notably in portable computers.\nRambus products sold on the market are organised in Rambus RIMM (trademark) memory modules, each module supporting 4, 6, 8, 12 or 16 Direct RDRAMs devices. RIMM modules are compatible with standard motherboard form factors; a motherboard usually supports up to three module sockets. The Direct Rambus Channel signals are daisy chained through each module. See Rambus RIMM Module Preliminary Information, document DL0078 available from Rambus Inc.\nThere is also provided in the Rambus specification a SPD (Serial Presence Detect) device. The purpose of the SPD is to store and provide sufficient information for a system to initialise the memory subsystem correctly: the SPD is a ROM device provided on each RIMM module, which includes information relating to the DRAM timing and device parameters, core organisation, module parameters, and other system level information. The SPD EEPROM devices of each RIMM module conform to the I2C wire protocol, and may be read into or written into by the memory controller of a Rambus system. See Direct Rambus SPD Specification 1.0, available from Rambus Inc.\nFIG. 2 is another view of a physical Rambus architecture, this time in an invalid configuration; it shows the memory controller 1 and the Direct Rambus Channel 2. Three modules 10-12 are connected to the bus; each side of each module may have up to 8 RDRAM devices, referenced again 3-6 on FIG. 2. Reference 13 is the SPD EEPROM of module 10; reference 14 shows the I2C protocol bus used by the memory controller for accessing the SPD EEPROMs of the different modules.\nMore details on Rambus may be found in the corresponding specification, issued by Rambus Inc. under the title Direct Rambus Technology Disclosure, Oct. 15, 1997.\nOne problem with Rambus is that the load of the bus is limited to 3 modules, and to 32 Direct RDRAM devices; if one of these limitations is violated, the bus system is not designed to be operational, and or even to boot at all; a computer in which the bus system is installed would in this case not be able to boot either.\nThis limitation on the number of modules is not likely in practice to be violated, since the bus normally comprises at most three module slots, and usually 2 or 3 module slots. However, a module may comprise up to 16 devices, so that the number of devices on the bus may exceed the highest allowable number of devices. This is the case in the configuration shown in FIG. 2.\nThus, the configuration of the bus hardware is such as to enable the bus to be improperly configured; in this case, a physical layer configuration constraint on the bus can be violated, and proper electrical operation of the bus is therefore not ensured. This can prevent the bus as a whole from booting properly. This possibility makes the system difficult for the user to upgrade or to diagnose problems that occur when they try to.\nA variety of bus configuration problems have been addressed in other contexts and a variety of solutions proposed.\nFor instance, U.S. Pat. No. 5,550,990 discusses physical partitioning of logically continuous buses. This document is directed to the SCSI (Small Computer System Interface) bus architecture, and suggests partitioning the bus into two or more physical entities which to the computer appear as one logical entity. This allows addressing problems potentially arising because of the scope of the architecture to be resolved; one example of such problems is excessive signal degradation due to use of signal rates which although allowed by the architecture are inappropriate for a particular bus loading. The solution disclosed in this document is to provide on the bus an adapter; instead of ensuring physical continuity of the bus, the adapter separates the bus into two bus partitions. This makes it possible, e. g. to operate the two partitions of the bus at different speeds, or to increase the number of devices connected to the bus. Where the speed has to be determined, a negotiation between the adapter and the devices connected to the bus is carried out at the time the adapter is initialised.\nU.S. Pat. No. 5,870,571 discusses automatic control of data transfer rates over a computer bus; this document is particularly directed to UltraSCSI buses. This document suggests detecting whether a SCSI external device is connected to the bus, and if this is the case, inhibiting the host adapter in order to reduce the data transfer rates to SCSI rate; otherwise, if no external SCSI device is detected, the UltraSCSI rate may be used, and the host adapter is not inhibited. In this document, the adapter polls the devices connected to the bus at initialisation, in order to know the transfer rate at which they may operate. Note that the operation of devices connected to the bus is not modified, since the host adapter only is inhibited.\nU.S. Pat. No. 5,237,690 discusses configuration at boot of IBM PS/2 personal computers. These computers provide a POS (programmable option select) for defining or providing settings for the assignment of system resources to a system board and various adapters. In order to avoid having to reconfigure the computer each time an adapter is added, removed or changed, this document suggests testing at boot of the computer whether any adapter was added, removed or changed; if this is the case, the adapters that were altered are disabled, and the computer is operated with all other adapters.\nU.S. Pat. No. 5,797,032 discusses a bus for connecting extension cards to a data processing system, and more particularly and ISA or EISA bus. For addressing the problem of collisions between the addresses of the different cards, this document suggests enabling all cards one at a time, for testing the addresses to which they respond. The cards that generate collisions are then disabled, and a message is displayed on a monitor for indicating to the user which cards were disabled.\nThe configuration constraints with which these two latter documents are concerned are logical-layer constraints and to resolve associated configuration problems the systems described rely on the buses concerned operating correctly at the physical level.\nAccording to the invention, there is provided a bus system comprising a controller; a high speed data transfer bus, the data transfer bus being subject to one or more inherent physical-layer configuration constraints for proper electrical operation; and a separate control bus, said control bus and said data transfer bus connecting the controller and the, or each, device connected thereto, wherein the controller is arranged to communicate with devices using the control bus in order to verify whether or not one or more of the physical-layer configuration constraints are satisfied and, if such configuration constraints are not satisfied, to modify using control signals transmitted on the control bus the operation of at least some of the devices in order to bring the data transfer bus to an operable condition.\nPreferably, if the configuration constraints are not satisfied, the controller is arranged to disable at least some of the devices using control signals transmitted on the control bus in order to bring the data transfer bus to an operable condition. The disabled devices may be the devices furthest from the controller on the data transfer bus. The controller may also be arranged to disable all devices connected to the bus, except one to five devices. The controller may also be arranged to set a stored indicator indicative of a error condition.\nIn one embodiment of the invention, the physical-layer constraints comprise a constraint on the number of devices connected to the bus.\nThe invention also provides a computer comprising such a bus.\nThe invention further relates to a process for bringing a data transfer bus to an operable condition in a bus system comprising a controller; a high speed data transfer bus, the data transfer bus being subject to one or more inherent physical-layer configuration constraints for proper electrical operation; and a separate control bus, said control bus and said data transfer bus connecting the controller and the, or each, device connected thereto. The process comprises the steps of\ncommunicating with devices using the control bus in order to verify whether or not one or more of the physical-layer configuration constraints are satisfied and,\nif such configuration constraints are not satisfied, to modifying the operation of at least some of the devices using control signals transmitted on the control bus.\nThe step of modifying may comprise disabling at least some of the devices using control signals transmitted on the control bus, and for instance, disabling devices furthest from the controller on the data transfer bus. The step of modifying may also comprises disabling all devices connect the bus, except one to five devices. The process may also comprise, if said configuration constraints are not satisfied, setting a stored indicator indicative of a error condition.\nIn one embodiment of the process, the physical-layer constraints comprise a constraint on the number of devices connected to the bus.\nThe invention also provides a computer program product for a computer with a bus system comprising a controller; a high speed data transfer bus, the data transfer bus being subject to one or more inherent physical-layer configuration constraints for proper electrical operation; and a separate control bus, said control bus and said data transfer bus connecting the controller and the, or each, device connected thereto. The computer program product comprises a computer readable medium having thereon:\ncomputer program code means, when said program is loaded, to make the controller communicate with devices using the control bus in order to verify whether or not one or more of the physical-layer configuration constraints are satisfied and,\nif such configuration constraints are not satisfied, to make the controller modify the operation of at least some of the devices using control signals transmitted on the control bus.\nPreferably, if such configuration constraints are not satisfied, the computer program code means make the controller disable at least some of the devices using control signals transmitted on the control bus. The disabled devices may be the devices furthest from the controller on the data transfer bus. The computer program code means may also make the controller disable all devices connected to the bus, except one to five devices. In another embodiment, the computer program code means set a stored indicator indicative of a error condition.\nThe physical layer constraints may comprise a constraint on the number of devices connected to the bus.\nThe invention provides a solution to the above described problem. The invention allows a computer at least to boot, even if the bus is improperly configured; this makes it possible to display a message to the user, so that he may address the problem. For Rambus, the mechanical configuration of the bus makes it possible to violate the bus specification by connecting an excessive number of devices on the bus.\nIn the case of the Rambus system, the limitation in the number of modules and devices connected to the bus is thought to be due to the sensitivity of the high speed signalling used on the Rambus Channel (the RSL or Rambus Signalling Levels) to the number of loads.\nIn consequence, a number of loads higher than the highest allowable number has no impact on the control bus, and does not affect the RSL signals for a few devices close to the controller. This makes it possible to disable some devices and to allow the bus to operate in a degraded operation mode. This mode is sufficient for booting a computer, and for allowing a warning to be communicated to the user, e. g. by displaying a message, so that they may reduce the number of devices on the bus and fix the problem.\nThe invention is however not limited to such a problem in the number of devices, but also can be applied in order to resolve other types of improper configuration; for instance, the invention could be applied if the bus can comprise different types of attachable devices, e.g. devices operating at different speeds, or devices requiring a special controller. It could also be applied for solving problems such as the mechanical length of a bus."}
-{"text": "The present invention relates generally to the field of can lids and specifically to lids for use on aluminum or metal beverage cans.\nThe applicant knows of no prior art which teaches the unique structure of his invention. U.S. Pat. No. 2,753,051 discloses a hinge on a pintle for sealing a container. Otherwise none of the structure of the present invention is shown. U.S. Pat. No. 3,372,832 (Yeater) discloses a removable cover for containers. Essentially it is a plastic cover having a pop down lid which locks into place by the force of the user pushing on the cap so that the projections 23 on the bottom side of the cap are pushed past the constricted middle section 16. The structure of the present invention is not disclosed. U.S. Pat. No. 4,331,255 (Fournier) discloses a two-part lid hinge. The two sections of this lid are connected by a hinge 16 which just may be a piece of plastic. The two sections are oriented so that the second section can be hooked over and on to the first so that the peripheral channel of the second section engages over the peripheral channel of the first section to provide a shallow space between the two sections bounded by a peripheral seal. Small openings are provided in a first section leaving a relatively large continuous imperforate area in that section. A removal tab is provided for forming a drinking opening in the second section, the tab being positioned so that it is disposed opposite the imperforate area in the first section. Again the structure of the present invention is not specifically disclosed. U.S. Pat. No. 3,977,559 (Lombardi) discloses a lid for a food container. It is a plastic cover that snaps over the top of a food container having a section of that lid that has been cut out so that it may be folded back along a hinge 22. The hinge does not have or disclose the present invention's structure nor does the cap itself disclose the structure of the present invention. U.S. Pat. No. 3,994,411 (Elfelt) discloses a lid for drink cups that includes a drinking flap of limited circumferential extent that may be selectively pivotally opened and closed. Such a drinking flap may be initially defined by frangible, i.e. breakable, lateral edges in the lid and may be held in its open position by the pull tab on the flap being inserted in a slit for a straw orifice. Essentially all this patent discloses, that is pertinent to the structure of the present invention, is a cap with hinges. The Elfelt structure does not appear to be re-usable. U.S. Pat. No. 4,284,200 (Bush) discloses a child resistant dispensing closure. Again, the structure of this invention is much different from the structure of the present invention and all that is really pertinent with this reference is the fact that there is a cap with a hinge. U.S. Pat. 4,361,250 (Foster) discloses a plastic container closure. This cap may be re-usable and has a flap or portion of the cover that flips open, similar in function to the present invention, in that the flap of Foster seals well. However, the structure of Foster is different from the structure of present invention. The hinged flap of Foster has hinge strips with depending pins formed along the sides of the flap that are integrally connected with the flap by tearable webs. After the flap has first been closed, the pins are anchored to the top of the closure and prevent the flap from opening during shipment of the container. Initial opening of the flap is effected by swinging the flap upwardly with a substantial force to tear the webs and separate the webs from the strips and the anchor pins. The torn webs provide visual indication that the flap has been opened. Accordingly, this is apparently a one use device or cap since once the webs have been torn the re-useability of the device is questionable. It is designed to be tamper resistant and tamper evident packaging. The structure of the present invention is not disclosed. U.S. Pat. 4,494,679 (Cleevely) discloses a thermoplastic container closure for dispensing solids. Again, this is another structure showing a hinged flap or flip cap. However, this flip cap does not appear to have a pintle hinge-type structure. U.S. Pat. No. 4,537,326 (Morehead) merely discloses a protector for a drink can opening. Specifically it is a device designed to attach to the portion of the lid where the pop tab or the pull tab is located so that it can be swung over or incorporated into the structure of the pop can lid in order to have a grating in place over the drinking opening and prevent the ability of insects like bees to enter into the container. Other than the slits the structure is completely different from the present invention. U.S. Pat. No. 4,582,214 (Dart) discloses a non-spill drink through lid. Slits to drink through are disclosed. No other structure of present invention is disclosed. U.S. Pat. No. 4,619,372 (McFarland) discloses a cap for a hot beverage cup. The cap is a disposable, removable closure cap for beverage containers and includes a depression permitting the beverage to be drunk while the cap remains in place on the container. The cap includes perforations in a depending wall located closely adjacent an inner wall of the container in order to limit the flow of beverage into the depression defined in the upper surface of the cap. A slit in the cap permits the aroma of the beverage to be enjoyed while the beverage is being drunk from the depression. The cap may be formed from sheet plastic material. This structure is completely different from the structure of the present invention, although it does disclose slits. It does not disclose the lip protection feature of the present structure although the well 28 of the McFarland device is designed to produce a somewhat similar effect. However the structure is completely different. U.S. Pat. No. 4,629,088 (Durgin) discloses a beverage container lid which includes a foldable flap which may be opened to allow the user to drink from a beverage container which is covered by the lid. A recess in the beverage lid is provided to receive the opened flap and to firmly secure the flap in its open position. The recess includes a pair of detentes on either side and an overhang at one end which cooperate to hold the flap firmly within the recess. The flap of course is also hinged. The structure of this cap is different than the structure of the cap which is the present invention, although it appears that recess 40 allows the flap to be flipped back and locked in place so that it is out of the way when a person drinks from the container. U.S. Pat. No. 4,949,865 (Turner) discloses a container lid with an integral stop. Essentially the structure of Turner is quite different from the present invention's structure and the main similarity is the fact that it discloses a hinged flap on the lid. The hinge is molded unitarily to the upper margin of the central support in the cover. Finally, U.S. Pat. No. 4,796,774 (Nabinger) also discloses a removable and re-sealable lid for a container but Nabinger's structure is much different from the present invention's structure."}
-{"text": "The polymerase chain reaction (PCR) is a method for increasing by many orders of magnitude the concentration of a specific nucleic acid sequence in a test sample. The PCR process is disclosed in U.S. Pat. Nos. 4,683,195; 4,683,202; and 4,965,188, each of which is incorporated herein by reference.\nIn PCR, a test sample believed to contain one or more targeted nucleic acid sequences is combined in a total volume of usually about 20 to 200 .mu.l with the following reagents: an aqueous buffer, pH 8-9 at room temperature, usually also containing approximately 0.05 M KCl; all four common nucleoside triphosphates (e.g., for DNA polymerase, the four common dNTPs: dATP, dTTP, dCTP, and dGTP) at concentrations of approximately 10.sup.-5 M to 10.sup.-3 M; a magnesium compound, usually MgCl.sub.2, usually at a concentration of about 1 to 5 mM; a polynucleotide polymerase, preferably a thermostable DNA polymerase, most preferably the DNA polymerase I from Thermus aquaticus (Taq polymerase and the Stoffel fragment of Taq polymerase are the subject of U.S. Pat. No. 4,889,818, incorporated herein by reference; the latter enzyme lacks the 5'.fwdarw.3' exonuclease activity of native Taq polymerase), usually at a concentration of 10.sup.-10 to 10.sup.-8 M; and single-stranded oligonucleotide primers, usually 15 to 30 nucleotides long and usually composed of deoxyribonucleotides, containing base sequences which are Watson-Crick complementary to sequences on both strands of the target nucleic acid sequence(s). Each primer usually is present at a concentration of 10.sup.-7 to 10.sup.-5 M; primers are synthesized by solid-phase methods well known in the art of nucleic acid chemistry.\nIn the simplest form, PCR requires two primers for each target sequence. These primers, when annealed to the opposing target strands, have their 3' ends directed toward one another's hybridization sites and separated by about 100 to 1,000 nucleotides (occasionally up to about 10,000 nucleotides). The polymerase catalyzes magnesium-dependent, template-directed extension of each primer from the 3' end of the primer, incorporating nucleoside monophosphates into the growing nucleic acid and releasing pyrophosphate.\nThis extension reaction continues until the polymerase reaches the 5' end of the template strand to which the extended primer was annealed, at which point the polymerase is free to bind to another primer-template duplex and catalyze extension of that primer molecule; the extension reaction also stops if the reaction mixture is heated to temperatures sufficient to separate the template from the extended primer before the enzyme has reached the 5' end of the template. After the enzyme has worked long enough to transform a large fraction of the primer-template duplexes into double-stranded nucleic acid, the latter can be denatured at high temperature, usually 90.degree. to 100.degree. C., to create two single-stranded polynucleotides, which, after cooling to a temperature where they can be annealed to new primer molecules, serve as templates for another round of enzyme-catalyzed primer extension. Because both DNA strands serve as template, each round of nucleic acid replication approximately doubles the concentration of the specific nucleic acid sequence defined at its ends by the two primer sequences. Therefore, the total concentration increase in the target nucleic acid sequence in a PCR amplification is by a factor of approximately 2.sup.n, where n is the number of completed thermal cycles between a high temperature where double-stranded DNA is denatured and a lower temperature or set of temperatures (40.degree. to 75.degree. C.) where primer-template annealing and primer extension occur.\nAlthough one can move PCR reaction tubes manually back and forth between thermostated baths in the two temperature ranges, PCR most commonly is performed in an automated temperature-controlled machine, known as a \"thermal cycler,\" in which a microprocessor is programmed to change the temperature of a heat-exchange block or bath containing reaction tubes back and forth among several specified temperatures for a specified number of cycles, holding at each temperature for a specified time, usually on the order of one-half to two minutes. Such a thermal cycler is commercially available from Perkin Elmer Cetus Instruments and described in the European Patent Publication No. 236,069 and U.S. patent application Ser. No. 670,545, filed Mar. 14, 1991, which is a continuation-in-part of Ser. No. 620,606, filed Nov. 29, 1990, both of which are incorporated herein by reference. The total cycle time is usually less than 10 minutes, and the total number of cycles is usually less than 40, so that a single, multi-cycle amplification, amplifying the targeted nucleic acid sequence 10.sup.5 to 10.sup.10 times, normally takes less than seven hours and often less than four hours.\nThe practical benefits of PCR nucleic acid amplification have been rapidly appreciated in the fields of genetics, molecular biology, cellular biology, clinical chemistry, forensic science, and analytical biochemistry, as described in the following review volumes and articles: Erlich (ed.), 1989, PCR Technology, Stockton Press (New York); Erlich et al. (eds.), 1989, Polymerase Chain Reaction, Cold Spring Harbor Press (Cold Spring Harbor, N.Y.); Innis et al., 1990, PCR Protocols, Academic Press (New York); and White et al, 1989, Trends in Genetics 5/6:185-189. PCR can replace a large fraction of molecular cloning and mutagenesis operations commonly performed in bacteria, having advantages of speed, simplicity, lower cost, and sometime increased safety. Furthermore, PCR permits the rapid and highly sensitive qualitative and even quantitative analysis of nucleic acid sequences, often with greatly increased safety because so much PCR product is made that nonisotopic detection modes suffice.\nDespite rapid and broad adoption of PCR by a range of biological and chemical disciplines, PCR has sometimes suffered from the occurrence of side reactions which interfere with amplification of the specific target sequence or sequences. Many amplifications yield non specific side products differing in size and sequence from the sequence targeted by the primers used. Sometimes nonspecificity is caused by mis-priming, where primers have been annealed to non-target sequences, also present in the nucleic acid of the test sample similar to the target sequence. Although the genomic DNA commonly contained in PCR test samples has customarily been thought to be completely double-stranded, procedures used to prepare DNA for amplification appear to render that DNA, to a significant extent, single-stranded. Single-stranded DNA is especially susceptible to mis-priming if mixed with a complete set of PCR reagents at ambient or sub-ambient temperatures. Many PCR reagents also yield primer dimers or oligomers, double-stranded side products containing the sequences of several molecules joined end-to-end, the yield of which correlates negatively with the yield of amplified target sequence.\nRecently several methodological modifications have improved PCR specificity and sensitivity significantly. In Hot Start.TM. PCR, complete mixing of PCR reagents and test sample is delayed until reactants have been heated to a temperature in the 50.degree. C.-80.degree. C. range, sufficient to minimize mis-priming and primer dimerization; thermal cycling is started immediately after mixing at elevated temperature. In manual Hot Start.TM. PCR, the operator heats the reaction tube, containing test sample and a subset of PCR reagents, to the elevated incubation temperature, opens each tube separately to add a small volume of liquid containing the missing reagent(s), and closes each tube before moving on to the next one. See Frohman et al., 1988, Proc. Natl. Acad. Sci. USA 85:8998-9002; Ward et al., 1989, Nature 341:544-546; Newton et al., 1989, Nucl. Acids Res. 17:2503-2516; and Faloona et al., Abstract 1019, 6th International Conference on AIDS, June 20-24, 1990, San Francisco, Calif. More recently, Hot Start.TM. PCR was rendered more convenient and precise by (1) replacement of the conventional mineral oil vapor barrier by a layer of wax melting in the 50.degree. C. to 80.degree. C. range, (2) assembly of reaction tubes such that before thermal cycling, PCR reactants are grouped into subsets separated by a solid wax layer, and (3) convective mixture of all reactants during the first heating step of thermal cycling after the solid wax melts into a lighter-than-water oil. Such wax-mediated, Hot Start.TM. PCR is the subject of U.S. patent application Ser. No. 481,501, filed Feb. 19, 1990, now abandoned in favor of continuation application U.S. Ser. No. 07/890,300, filed May 27, 1992, incorporated herein by reference.\nAlternatively, nonspecific amplified nucleic acid resulting from primer dimerization and mis-priming while completely mixed PCR reactants stand at room temperature before thermal cycling can be destroyed by an enzymatic restriction process described in PCT U.S. patent application Ser. No. 91/05210 filed Jul. 23, 1991, which published as PCT Patent Publication No. WO 92/01814 on Feb. 6, 1992, which is a continuation-in-part of U.S. Ser. No. 609,157, filed Nov. 2, 1990, now abandoned which is a continuation-in-part of U.S. Ser. No. 557,517, filed Jul. 24, 1990, now abandoned each of which is incorporated herein by reference. To perform such restriction, one of the conventional four dNTPs is replaced by a structural analogue which is incorporated into all amplified nucleic acid by the PCR polymerase. Also included in the reaction mixture is an enzyme which digests nucleic acid at (and only at) positions which contain the structural analogue; this enzyme must be active only at temperatures below about 50.degree. C., so that it does not damage amplified nucleic acid during thermal cycling at higher temperatures. Preferably the restriction enzyme is permanently inactivated during thermal cycling, so that it cannot damage amplified nucleic acid if the latter is stored for any significant period of time at room temperature after amplification and before analysis. The most practical restriction enzymes are glycosidases which cleave from the polynucleotide phosphodiester backbone the unconventional nucleic acid base introduced by the dNTP analogue. The resulting abasic sites experience cleavage of the polynucleotide phosphodiester backbone upon heating. This restriction process has been integrated practically with PCR by replacing dTTP with dUTP and by incorporating in the reaction mixture the enzyme uracil-N-glycosidase.\nA chemical variant of the Hot Start.TM. process incorporates into the PCR reagent mixture a single-stranded DNA binding protein (SSB) at a concentration sufficient to bind a significant fraction of the single-stranded DNA present before thermal cycling is started. This ssDNA comprises minimally the primers, the concentrations of which are well known by the operator, and may also include slight or considerable amounts of the test sample DNA, depending on whether the latter has been prepared in a way which might denature it. During thermal cycling, the binding of the SSB to primers and single-stranded template strands formed by PCR product denaturation must be weak enough not to interfere with primer-template annealing and enzymatic primer extension. Before thermal cycling, while reactants stand together at room temperature, SSB binding to the primers and any single-stranded regions of test sample DNA must be strong enough to block mis-priming and primer dimerization. Two heavily studied SSBs (Chase and Williams, 1986, Ann. Rev. Biochem. 55:103-136) are commercially available and have been used with PCR: gene 32 protein from the bacteriophage T4 and the 19 kilodalton SSB from E. coli (19 kda is the subunit size; the normal active species is a tetramer). SSB is the major active ingredient of Perfect Match.TM. polymerase enhancer, a mixture of E. coli SSB and bovine serum albumin sold by Stratagene (San Diego, Calif.) for the purpose of increasing PCR specificity and yield. Bacteriophage gene 32 protein has been included in PCR mixtures to improve amplification of long targets (Schwarz et al., 1990, Nucl. Acids Res. 18:1079) and to relieve polymerase inhibition by blood in the test sample (Panaccio and Lew, 1991, Nucl. Acids Res. 19:1151). However, essentially all organisms possess SSBs with compositions unique to each organism. Other SSBs which have been characterized biochemically include one from a filamentous bacteriophage (Brayer and McPherson, 1984, Biochemistry. 23:340-349), a family of sequence-homologous proteins from plant virus (Saito et al., 1988, Virology 167:653-656, and Citovsky etal., 1990, Cell 60:637-647), and one from Agrobacterium tumefaciens (Citovsky et al., 1989, Proc. Natl. Acad. Sci. USA 86:1193-1197). SSBs possess enough structural similarity to suggest that DNA binding is associated with a consensus structure of alternating aromatic amino acids (phenylalanine, tyrosine, and tryptophan) and charged amino acids (glutamate, aspartate, lysine, and arginine) (Prasad and Chiu, 1987, J. Mol. Biol. 193:579-584) such that artificial polypeptides might be created which function as well as the biological SSBs in improving PCR specificity and yield. In addition, enough is known about SSB structure and function to suggest ways to improve function by genetic engineering.\nAlthough the three basic tactics of PCR specificity enhancement (Hot Start.TM. methods, amplified DNA restriction, and SSB addition to the reaction mixture) each can serve alone to improve specific amplification, combinations of the three approaches may have special benefits. For example, whereas Hot Start.TM. methods block only that nonspecificity resulting from reactant incubation at ambient temperature before cycling is started, SSB s may reduce mis-priming which arises during thermal cycling. On the other hand, SSB used without a manual or wax-mediated Hot Start.TM. process occasionally will trigger massive primer dimerization which interferes with specific amplification. The combination of the two methods optimally reduces mis-priming and primer dimerization.\nThe preceding background art has dealt with conventional PCR, wherein test sample nucleic acids are extracted from a biological source in a way which destroys target sequence association with individual cells or subcellular structures. So-called in situ nucleic acid hybridization methods have evolved to detect target sequences in the cells or organelles where they originated (for a review of the field, see Nagai et 1987, Intl. J. Gyn. Path. 6:366-379). Typically, in situ hybridization entails (1) preparation of a histochemical section or cytochemical smear, chemically fixed (e.g., with formaldehyde) to stabilize proteinaceous subcellular structures and attached to a microscope slide, (2) chemical denaturation of the nucleic acid in the cellular preparation, (3) annealing of a tagged nucleic acid probe to a complementary target sequence in the denatured cellular DNA, and (4) localized detection of the tag annealed to target, usually by microscopic examination of immobilized nonisotopic (absorbance or fluorescence staining) or isotopic (autoradiographic) signals directly or indirectly generated by the probe tag. However, conventional in situ hybridization is not very sensitive, generally requiring tens to hundreds of copies of the target nucleic acid per cell in order to score the presence of target sequence in that cell.\nRecently, the sensitivity enhancement associated with PCR amplification of target sequence has been combined with the target localization of in situ hybridization to create in situ PCR, wherein PCR is performed within chemically fixed cells, before (Haase et al., 1990, Proc. Natl. Acad. Sci. USA 87:4971-4975, incorporated herein by reference) or after (Nuovo et al., 1991, Amer. J. Pathol. in press, incorporated herein by reference) the fixed cells have been attached to a microscope slide; the amplified nucleic acid is located by microscopic examination of autoradiographs following isotopically tagged probing (Haase et al., supra) or stained patterns directly deposited on the microscope slide following enzyme-linked detection of biotin-tagged probes (Nuovo et al., supra). The cells may be suspended (Haase et al., Supra) or may be part of a tissue section (Nuovo et al., supra) during in situ amplification.\nIn situ PCR requires a delicate balance between two opposite requirements of PCR in a cellular preparation: the cell and subcellular (e.g., nuclear) membranes must be permeabilized sufficiently to allow externally applied PCR reagents to reach the target nucleic acid, yet must remain sufficiently intact and nonporous to retard diffusion of amplified nucleic acid out of the cells or subcellular compartments where it is made. In addition, the amplified nucleic acid must be sufficiently concentrated within its compartment to give a microscopically visible signal, yet remain sufficiently dilute that it does not reanneal between the denaturation and probe-annealing steps. Haase et al., supra, relied on paraformaldehyde fixation of cells to have created sufficient but not excessive permeability. Nuovo et al., supra, also employed a single, commercially available, proteinase treatment to improve permeability.\nBoth Haase et al., supra, and Nuovo et al., supra, used a series of PCR primer pairs to specify a series of overlapping target sequences within the genome of the targeted organism to improve retention of amplified target nucleic acid within the cellular compartment where it was made. The resulting PCR product was expected to be so large (greater than 1,000 base pairs) that its diffusion from site of origin should be greatly retarded. However, the use of multiple primer pairs severely reduces the practicality of in situ PCR, not just because of the expense associated with producing so many synthetic oligonucleotides, but even more seriously because many PCR target organisms, especially pathogenic virus, are so genetically plastic that it is hard to find even a few short sequences which are sufficiently invariant to make good primer and probe sites. Other important target sequences, such as activated oncogenes, inactivated tumor suppressor genes, and oncogenic chromosomal translocations, involve somatic point mutations and chromosomal rearrangements which can be distinguished from the parental sequence if relatively short PCR products are amplified from single primer pairs. Multiple primer pairs and long structures would frustrate attainment of the specificity often needed to distinguish cancerous cells from their normal neighbors. Multiple primer pairs jeopardize PCR in a different way as well; they promote primer dimerization and mis-priming, reducing sensitivity and specificity and increasing the likelihood of false-negative results because nonspecific amplification radically reduces the yield of amplified target sequence. Reinforcing the tendency of multiple primer pairs to enhance nonspecific amplification are the rather high primer concentrations preferred for in situ PCR (Nuovo et al., supra).\nOne useful variant of conventional PCR detects target RNA sequences in test samples by creating complementary DNA (cDNA) sequences with the catalytic mediation of added reverse transcriptase; the cDNA then is subjected to standard PCR amplification (Kawasaki et al., 1988, Proc. Natl. Acad. Sci. USA 85(15):5698, and Rappolee et al., 1989, J. Cell. Biochem. 39:1-11). Recently, such RNA PCR has been streamlined by using a thermostable DNA polymerase which, depending on exact chemical conditions, also shows strong reverse transcriptase activity. This enzyme and its optimized application to RNA PCR are subject of PCT U.S. patent application Ser. No. US90/07641, filed Dec. 21, 1990, incorporated herein by reference. Adaptation of in situ PCR to RNA targets will realize the full potential of the method to differentiate among neighboring cells in a histochemical or cytochemical preparation with respect to somatic mutation, pathogenic infection, oncogenic transformation, immune competence and specificity, state of differentiation, developmental origin, genetic mosaicism, cytokine expression, and other characteristics useful for understanding both normal and pathological conditions in eukaryotic organisms.\nThe present invention increases the convenience, sensitivity, and specificity of in situ PCR, also eliminating any need for multiple primer pairs to detect a single target sequence. In doing so, it also allows in situ PCR to discriminate among alleles and increases the practicality of in situ PCR analysis of RNA targets."}
-{"text": "Various devices and systems are known for use in harvesting solar energy by the use of photovoltaic cells. These include slat concentrators, which are photovoltaic devices generally comprising a series of parallel trough-shaped off axis parabolic reflectors to concentrate sunlight on photovoltaic receptors mounted on respective adjacent reflectors. The reflectors are typically automatically actuated to track the sun in order to ensure accurate reflection and concentration of solar radiation on the photovoltaic receptors.\nThe photovoltaic receptors forming part of such concentrators have a limited lifespan and the photovoltaic devices therefore require periodic removal and replacement of the photovoltaic receptors. There is a relationship between the operating temperatures of the photovoltaic receptors and their lifespan. Additionally, a photovoltaic receptor generally displays higher efficiency at lower temperatures."}
-{"text": "Online auctions delivered to users via the World Wide Web have become a popular activity for users interested in purchasing goods and services at the least possible amount online. Many online auctions allow a person to bid on an item in an attempt to obtain a winning bid and an associated opportunity to purchase the auction item at the winning bid price. Known online auction websites allow several users to compete with other online auction participants by submitting bids on a particular item until a predefined period of time has elapsed, and the highest bidder is determined.\nAt least some known online auction websites require users to purchase a bid unit that represents an opportunity to submit a bid in an auction and to redeem a bid unit each time the user places a bid on an auction item. Once the user has redeemed all of the bid units, the user cannot participate in the auction until additional bid units are purchased or provided. By requiring a user to purchase bid units to participate in the auction the online auction receives revenue from each user participating in the auction.\nMany auction participants are attracted by the opportunity to win new auction items of high value at a low price, but may become frustrated with a requirement to purchase new bid units each time a user wishes to place a bid on an auction item. Auction participants may also become frustrated by purchasing bid units, participating in the auction, and not winning an auction. Likewise, online auction providers desire new auctioning opportunities to appeal to their auction participants, entice users to place additional bids, spend more time in the auction website, and have additional opportunities to benefit from winning new auctions. Accordingly, there is a continued need for systems and methods that create, provide, and facilitate new and interesting online auctions that are fair to each customer.\nSome online auction formats, such as the penny auction, also known as a click to bid auctions have been known to be subject to fraud, cheating, and abuse by both customers and auction operators. The two most common forms of cheating are known as \u201cshill\u201d bidders and \u201cbot\u201d bidders. A \u201cshill\u201d bidder is used by an auction operator to artificially inflate the price of the auction while increasing the amount of the customer bids by having employees of the auction operator pose as bidders and bid against the real customers. A \u201cbot\u201d bidder accomplishes the same result, however the false bidding is accomplished by software code, or script, which poses as a real person competitively bidding in the auction. The reputation of the penny auction industry has been severely damaged by the use of shill and bot bidders. The reputation of fraud in the industry has been so damaging that many people refuse to participate in a penny auction for fear of being cheated out of their money.\nAnother detractor to acquiring new customers and retaining existing penny auction customers is the potential for any customer to be out bid by a well-financed competitor customer, thereby winning most auctions. For example, a customer who has the personal financial ability to purchase ten times the bids of a lesser well-financed customer is at a distinct and nearly insurmountable advantage. The better financed customer could simply keep bidding in an auction against the lesser financed customer until he or she runs out of bids and is unable to purchase more bids to remain competing in the auction. This appearance that penny auctions unfairly provide advantage to well-funded customers is known to prevent many new potential customers from trying the product for a first time, thereby severely limiting the potential growth of a penny auction business.\nThe present invention is aimed at one or more of the problems identified above."}
-{"text": "The present invention relates to a bumper system for motor vehicles.\nAs is known, the main function of bumper systems for motor vehicles is to convert impact energy to deformation energy in a collision with another vehicle or with any stationary or non-stationary obstacle. This limits damage to the motor vehicle.\nAs a rule, one known embodiment of such a bumper system comprises a cross-member that has a U-shaped cross-section, with the bar portion thereof connected via crush boxes to the longitudinal supports of the motor vehicle. The ends of the leg portions of the cross-member are connected, especially joined, using a closure plate, such that the closure plate prevents the cross-member from opening and/or flattening under a bending load. The profile of the bumper system is closed and thereby achieves high flexural rigidity, while remaining lightweight.\nBoth steel alloys and aluminum alloys are suitable materials for the cross-member and the closure plate. Furthermore, combinations of these materials are possible."}
-{"text": "Air filtration masks (referred to herein as \u201cfilter masks\u201d) are widely used to protect people from air borne contaminants and gasses. For example, air borne dust particles are a known hazard commonly on work sites. Consequently, workers normally wear filter masks to avoid inhaling the dust particles. To that end, filter masks used in this application are manufactured with a filter material specified to prevent, among other things, a substantial majority of dust particles from being inhaled by the worker.\nIn addition to primarily filtering inhaled air, some filter masks are specifically manufactured to filter both inhaled and exhaled air. For example, hospital staff often wear filter masks to prevent both their germs from infecting patients, and patients' germs from infecting them.\nThere is a need in the art to improve the filtration efficiency of filter masks. Accordingly, filter masks with a higher efficiency filter layer and/or multiple filter layers have been developed for that purpose. However, this often has the undesirable effect of increasing the air resistance through the filter mask and may cause several problems.\nFor example, a person wearing the filter mask may have a more difficult time breathing due to the increased air resistance. To overcome this problem while still providing improved filtration efficiency, filter masks have been developed that have an increased filter area. Manufacture of such filter masks, however, can be quite complex. For example, increasing the filter area can cause various portions of the filter layer to overlap or can be costly to construct. Overlap can effectively increase the thickness of the filter layer, thus causing the same air resistance problem as discussed above.\nAdditionally, since a person wearing the mask while performing manual labor must typically breathe heavier, the filter layer(s) is more likely to flex and eventually collapse around the face. This collapse may cause portions of the face mask to contact and irritate the face of the person wearing the face mask, as well as cause discomfort. Consequently, efforts have been made to stiffen the mask, such as by adding additional material to the filter mask. However, adding additional material to the face mask adds complexity to the production process and increases cost."}
-{"text": "Field of the Invention\nThe present invention generally relates to a vehicle door structure. More specifically, the present invention relates to vehicle door structure having a first bracket and a reinforcing bar with a contoured end surface that is spaced apart from adjacent surfaces of the first bracket in an undeformed state, the contoured end surface being shaped to maximize surface area contact between the first end section of the reinforcing bar and the adjacent surfaces of the first bracket in response to a force being applied to the reinforcing bar and causing the first bracket to deform.\nBackground Information\nMany areas of a vehicle body structure are provided with reinforcing structures. For example, vehicle doors usually include reinforcing structures."}
-{"text": "The instant invention relates to exercise apparatus for exercising the adductor muscles of the legs and more particularly to an adductor exercise apparatus having means for adjusting the angular starting position of the leg receiving members.\nExercise machines, and more particularly, adductor exercise machines, have heretofore been known in the art. In this connection, the U.S. Patents to Scott No. 4,022,463; DeNiro No. 4,892,304; Dela Rosa No. 4,877,239; and Goodman No. 5,026,049 represent the closest prior art to the subject invention of which the applicant is aware.\nThe patent to Scott discloses an exercise apparatus for exercising the adductor muscles of the legs, however it does not disclose any means for adjusting the starting position of the leg receiving means. The patent to DeNiro discloses an exercise apparatus including a plurality of movable parts which are adjusted by means of pin and hole mechanisms. The patent to Dela Rosa discloses a stretching apparatus having an elongated threaded adjusting shaft in the form of a worm gear for adjusting the degree of split of the leg receiving elements. The patent to Goodman shows leg receiving pads which are adjustable to accommodate different size persons.\nThe instant invention provides an adductor exercise apparatus which includes an adjustment mechanism for adjusting the angular starting positions of the leg receiving assemblies. More specifically, the adductor exercise apparatus includes a base, two support legs which extend outwardly from opposite sides of the base, a seat on the base, and two leg receiving assemblies for receiving the legs of a user seated on the seat. The leg assemblies are pivotably mounted to the base so that they are pivotably movable between a spread apart position and a parallel, together position. The leg assemblies include a pad assembly including two pads which respectively engage with the thigh and calf portions of the user's leg. The calf pad is slidably movable with respect to the thigh pad. The apparatus further includes a cabled weight assembly for normally urging the leg assemblies toward the spread apart position. The adjustment mechanism consists of a pair of movable adjustment arms each having a first end which engages the corresponding leg receiving assembly, and a second end which is slidably received and secured in a sleeve assembly mounted on the support leg adjacent to the respective leg receiving assembly. The arms are slidably adjustable in the sleeve assemblies to a plurality of predetermined longitudinal positions. In this manner, as the leg receiving assemblies are urged toward their normal spread apart position, the leg receiving assemblies engage with the projecting ends of the adjustment arms to position the leg receiving assemblies in corresponding angular starting positions. The cabled weight assembly includes stacked weight members and a cable system which extends around a series of pulleys for translating pivoting movement of the leg receiving assemblies into vertical movement of the weight members. The weight members are lifted by an elongated weight bar which is attached to the cable system.\nAccordingly, it is an object of the instant invention to provide an exercise apparatus for exercising the adductor muscles of the leg.\nIt is another object to provide an adductor exercise apparatus having an adjustment mechanism for adjusting the angular starting positions of the leg receiving assemblies.\nIt is yet another object to provide an adductor exercise apparatus including leg receiving means having two separate pads for engaging with the thigh and calf portions of the user's leg.\nIt is still another object to provide an adductor exercise apparatus having a cabled weight assembly and an elongated weight bar for lifting the weight members of the weight assembly.\nOther objects, features and advantages of the invention shall become apparent as the description thereof proceeds when considered in connection with the accompanying illustrative drawings."}
-{"text": "1. Field of the Invention\nThe present invention relates to a press for the stamping of plates, particularly vehicle license plates having license numbers. The license plates are made from sheet metal material, itself in the form of strip material or individual plates, that is fed in synchronized manner into the press. The press includes a press stand comprised of a press table and a press stamp, as well as interchangeable stamping tools.\n2. Description of the Related Art\nPresses of the above type, as found in DE 32 03 801 C2, are furnished with tools, embodied as block tools, having a uniform width. The tools are used for stamping license numbers (having letters and numbers) in vehicle license plates and do so in registered layout, or in a registered print style, as known in Germany.\nIn other countries such as Austria, the figures on the license plate are stamped in a non-registered layout, or in a free print style. In other words, the stamping tools have a width that can be tailored to the width of letters, numbers and, if need be, coats of arms.\nThe present invention seeks to develop a press for the stamping of plates, particularly vehicle license plates with license numbers in any linear order, so that it can be furnished with stamping tools of uniform width (i.e. uniform registration) as well as with tools having different widths."}
-{"text": "1. Field of the Invention\nExemplary embodiments relate to instant messages, and particularly to redirecting instant messages from one system to another system.\n2. Description of Related Art\nThere exists a growing popularity in instant messaging services. Instant messaging (IM) is a form of real-time communication between two or more people based on typed text. The text is conveyed via computers connected over a network such as the Internet.\nInstant messaging offers real-time and/or near-time communication and allows easy collaboration, which might be considered more akin to genuine conversation than email's \u201cletter\u201d format. In contrast to e-mail, the parties know whether the peer is available because most systems allow the user to set an online status or away message so that peers are notified when the user is available, busy, or away from the computer.\nInstant messaging allows instantaneous communication between a number of parties simultaneously, by transmitting information quickly and efficiently, featuring immediate receipt of acknowledgment or reply. In certain cases IM involves additional features, which make it even more popular, i.e., to see the other party, e.g. by using web-cams, or to talk directly for free over the Internet.\nCurrently, when utilizing instant messaging from one computer and then subsequently logging into the instant messaging session from another computer, the first computer may automatically be logged off. Although new messages may be sent to the second computer, messages already delivered to the first computer would remain on the first computer screen until the user returns to the first computer and dismisses them. Thus, the delivered messages would not be delivered to the user in a timely fashion.\nSome messengers permit being logged on from multiple (locations) computers at once and allow the current chat record to be fully available at all locations. This requires the instant messenger to support synchronous logins from multiple devices, however, which may be unwanted complexity to support, or may be considered a security issue. Without the extra complexity and security issues, it would be desirable to have a method to prevent messages from being untimely even if messages have been delivered to a computer that has been automatically logged off."}
-{"text": "A non-aqueous electrochemical cell using lithium as an active anode material has high energy density, good storage characteristics and wide operation temperature range. A non-aqueous electrochemical cell is therefore often used as a power source for a calculator, a watch or a memory back up system. Such a cell comprises an anode, an electrolyte and a cathode. In general, such a cell uses as an anode an alkali metal such as lithium or sodium; as an electrolyte or electrolytic solution, a solution of a solute such as lithium perchlorate or lithium tetrafluoroborate in a non-aqueous solvent such as propylene carbonate, .gamma.-butyrolactone or diglyme; and as a cathode, manganese dioxide or polycarbonmonofluoride.\nThe combination of relatively high theoretical energy density, potentially long life, and low cost materials such as reported in the sodium-sulfur system high temperature batteries is suitable primarily for low rate performance work such as electric road vehicle propulsion or load leveling of electric power supplies. The sodium-sulfur systems, first proposed in 1966, have had a great deal of effort expended in trying to develop a practical system. The basic operating principle involves the separation of two active molten materials, sodium and sulfur, by either a ceramic membrane of beta alumina or sodium glass, which at about 300.degree. C. or higher allows the passage of sodium ions that form with the sulfur any of the several polysulfides. The open circuit voltage of the system is at just over 2 volts, about the same as the lead-acid cell. Two formidable problems exist at the present time, viz., cracking of the separator and corrosion of the casing and seal.\nAnother somewhat similar system is the lithium-iron sulfide system, operating at about 450.degree. C. However, insufficient development has been done to date to demonstrate the widespread practicality of this system.\nAnother of the developments being pursued involves a lithium-based cell, in which the negative electrode is a lithium alloy (typically either lithium-aluminum or lithium-silicon), the positive electrode is an iron sulfide, and the electrolyte is a molten salt, such as the eutectic composition in the lithium chloride potassium chloride system. Because of the high melting point of such salts, such cells must be operated in the temperature range of 400-500 degrees centigrade.\nThis requirement to operate at such high temperatures has several important disadvantages. One of these is that various degradation processes, such as corrosion of the cell container, seals, and other components are accelerated by such high temperatures. Another is that a substantial amount of energy is lost through heat transfer to the surroundings. Still another is that the voltage obtained from such cells is lower at elevated temperatures, due to the fundamental property of the negative temperature dependence of the free energy of the cell reaction. Furthermore, the higher the temperature of operation, the greater the potential problems related to damage to the cell during cooling to ambient temperature and reheating, whether deliberate or inadvertent. Differences in thermal expansion, as well as dimensional changes accompanying phase changes, such as the freezing of the molten salt, can cause severe mechanical distortions, and therefore damage the cell components.\nCells involving a lower temperature molten salt electrolyte have been investigated where the molten salt is based upon a solution of aluminum chloride and an alkali metal chloride. However, the soluble positive electrode materials are not stable in the presence of the respective alkali metals. As a result, an auxiliary solid electrolyte must be used to separate the alkali metal and the counter electrode. One example of such a cell involves a molten sodium negative electrode, a solid electrolyte of sodium beta alumina, a molten aluminum chloride-sodium chloride salt, and either antimony chloride or an oxychloride dissolved in the chloride salt as the positive electrode reactant. Such a cell can operate in the temperature range of 150-250 degrees centigrade. It has the disadvantage of having to employ an electrolyte, which increases the cell impedance, as well as adding to the cost and complexity.\nU.S. Pat. No. 3,844,837 to Bennion et al discloses a non-aqueous battery in which the anode may be lithium and/or graphite on which lithium metal is deposited and as a positive electrode a platinum cup filled with powdered K.sub.2 SO.sub.4 and graphite is utilized. The electrolytes disclosed are LiClO.sub.4, LiCF.sub.3 SO.sub.3 or LiB.sub.4 dissolved in dimethyl sulfite.\nU.S. Pat. No. 4,877,695 discloses a non-aqueous electrochemical cell for use in primary rechargeable storage devices in which the cathode comprises an electrically conductive carbonaceous material, the anode is a molten mixture of two elements selected from the group consisting of sodium, potassium, cesium and rubidium, and the electrolyte comprises a solvent and an electrolyte salt selected from the group consisting of an alkali metal tetrafluoroborate and a tetraalkylammonium tetrafluoroborate. The present invention provides a specific improvement in the electrochemical cell of the patent by increasing its energy and power density.\nU.S. Pat. No. 4,886,715 discloses a primary rechargeable energy storage device comprising an anode which is an alkaline earth or alkali metal and a carbonaceous fiber cathode. The electrolyte is a membrane comprising lithium laurate."}
-{"text": "The delivery of text-based messages (i.e. data packets) from a sending device to one or more receiving devices over a wireless LAN, presents special challenges. Typically, the message is routed through a wireless gateway where it is stored until it has been transmitted to, and stored within, an electronic mail server of the data network. Receiving devices are then able to retrieve stored messages from the electronic mail server at their convenience. The speed at which electronic messages are transmitted from a sending to a receiving device depends in part on how efficiently data packets are transported from a sending mobile device to an electronic mail server through wireless communication networks.\nWhen a wireless gateway receives a data packet from a mobile device over a wireless network, the received data packet is sent to a destination electronic mail server. However, to ensure that the data packet is successfully transmitted to the destination electronic mail server and not lost in the meantime, the wireless gateway generally stores the data packet in an internal permanent storage device (e.g. a database server or a file system) before transmitting the data packet. Typically, the wireless gateway waits until the permanent storage device confirms storage of the data packet before proceeding with processing the data packet or even with sending acknowledgement of the receipt of the data packet back to the mobile device. This kind of storage procedure appreciably slows down the processing of data packets within the router."}
-{"text": "In order to restore or assume the function of a damaged or lacking tissue or organ the transplantation of natural tissues or organs from another donor or, if possible, the concerned individual himself is an established practice that has long been known. Due to the constant lack of suitable donor tissues and donor organs and other disadvantages of natural tissue, e.g., rejection reactions and the risk of the transfer of diseases from the donor to the recipient, many efforts are being directed to the production of artificial tissues and organs as alternatives.\nA customary procedure for this is the provision of a carrier or matrix structure from material compatible with the body that is colonized with differentiated cells of the target tissue, and the cultivation of the cells in vitro until a tissue-like cell structure has been produced. The differentiated cells are obtained either from cultures of explanted tissue samples or from stem cells that had been stimulated to differentiate. The use of stem cells permits a more rapid production of larger amounts of the desired cells in many instances. Traditionally, pure cell populations of a certain type are produced. Most of the in vitro organs or in vitro tissues known in the state of the art are disadvantageous in as far as that they do not have or do not develop the tissue structure that corresponds to the morphological constitution of the native tissue or organ even after implantation and a fairly long residence time in the body. This applies as a rule even when the carrier matrix had been colonized with several different populations of tissue-typical cells.\nAnother approach for treating degenerative diseases or damage to tissues and organs using stem cells consists in implanting the stem cells and/or differentiated cells derived from them directly into the damaged tissue/organ in order to proliferate there and result in a regeneration of the damaged tissue/organ. In this instance too the implantation or regeneration of a certain cell type traditionally is in the foreground. The problem is, however, that in very many degeneration phenomena or damages to various organs a plurality of cells is always involved so that, e.g., in the skin keratinocytes, epithelial cells and blood vessels are also affected in addition to fibroblasts. In the case of nerve damage even glia cells often have to be replaced too and muscle damage often means destruction of the accompanying nerves."}
-{"text": "A superacid is a known class of acidic material that has an acid strength, measured by Hammett acidity function H.sub.0, greater than that of 100% H.sub.2 SO.sub.4, which has an H.sub.0 value of -12. A superacid, therefore, has an H.sub.0 value of less than -12 or an acid strength greater than -12. Superacids are useful for reactions that are generally catalyzed by strong acid, such as paraffin isomerization.\nM. Hino and K. Arata describe the synthesis of a solid superacid having an acid strength of up to H.sub.0 .ltoreq.-16.04 by exposing hydroxides or oxides of Fe, Ti, Zr and Hf, prior to crystallization, to sulfate ions, followed by calcination in air at over 500.degree. C. in J. Chem. Soc., Chem. Commun., 1148 (1979). K. Arata and M. Hino also describe the synthesis of a solid superacid having an acid strength of H.sub.0 .ltoreq.-14.52 by impregnating Zr(OH).sub.4 or amorphous ZrO.sub.2 with aqueous ammonium metatungstate, followed by calcining in air at 800.degree. to 850.degree. C. in J. Chem. Soc., Chem. Commun., 1259 (1988) and in \"Synthesis of Solid Superacid of Tungsten Oxide Supported on Zirconia and Its Catalytic Action\", Proceedings 9th International Congress on Catalysis, Volume 4, pages 1727-1734 (1988).\nThe superacid described by K. Arata and M. Hino in \"Synthesis of Solid Superacid of Tungsten Oxide Supported on Zirconia and Its Catalytic Action\" is particularly useful as a catalyst in the isomerization of butane to isobutane and of pentane to isopentane. However, in order to obtain maximum catalytic activity, K. Arata and M. Hino report calcination temperatures of from 800.degree. to 850.degree. C. for the tungsten-modified zirconia catalyst.\nSoled et al describe, in U.S. Pat. No. 5,113,034, a solid acid catalyst having an H.sub.0 value ranging from -14.5 to about -16.5, comprising a sulfate or tungstate-modified Group IVB oxide, which is useful to dimerize C.sub.3 or C.sub.4 containing feedstreams. A calcination temperature range for the tungstate-modified zirconia catalyst of 450.degree. to 800.degree. C. is given. Specifically, the Examples use an initial calcination temperature of 600.degree. C., followed by a calcination temperature of 800.degree. C. prior to charging.\nU.S. Pat. No. 5,198,403 (Brand et al, herinafter Brand) discloses a catalyst for the selective reduction of nitrous oxide with ammonia which contains, in addition to titanium oxide as component A, at least one oxide of tungsten, silicon, boron, aluminum, phosphorus, zerconium, barium, yttrium, lanthanum, cerium as component B1 and at least one oxide of vanadium, niobium, molybdenum, iron and copper as B2, with an atomic ratio between the elements of components A and B from 1:0.001 up to 1. The catalyst of Brand requires the presence of titanium oxides, presumably because it is effective in the reduction of nitrous oxides. Reduction of nitrous oxides is the intended purpose of Brand. The catalysts of the instant invention do not contain titanium and are not drawn to nitrous oxide reduction.\nU.S. Pat. No. 4,918,041 (Hollstein et al, hereinafter Hollstein) is directed to a sulfated calcined solid catalyst. The catalysts of the instant invention are not sulfated. Group VI B metals such as tungsten are employed in the superacid, rather than Group VI A metals such as sulfur. Group VI B metals are preferred in the instant invention due to their stability and ease in regeneration.\nThe generation of acid activity in solid oxide catalysts in general, and in the tungsten-modified zirconia catalyst specifically, requires calcination of the catalyst at temperatures of about 800.degree. C. This extreme temperature, however, causes significant loss of catalyst surface area. For example, K. Arata and M. Hino report surface areas of 35.3 and 29.5 m.sup.2 /g for catalysts calcined at 800.degree. and 900.degree. C. in \"Synthesis of Solid Superacid of Tungsten Oxide Supported on Zirconia and Its Catalytic Action\". By contrast, the surface areas for catalysts calcined at 600.degree. and 700.degree. C. are reported as 44.2 and 38.5 m.sup.2 /g. Hence, it is clear that as the calcination temperature is increased, the surface area of the catalyst is decreased. Further, the extreme calcination temperature required to generate acid activity results in a more difficult manufacturing process."}
-{"text": "This application is based upon and claims the benefit of priority from the prior Japanese Patent Applications No. 2000-117739, Apr. 19, 2000; No. 2000-117740, Apr. 19, 2000; No. 2001-055815, Feb. 28, 2001; and No. 2001-110469 Apr. 9, 2001, the entire contents of which are incorporated herein by reference.\n1. Field of the Invention\nThe present invention relates to a focus stabilizing apparatus for stabilizing a focused state of an optical apparatus, such as a microscope.\n2. Description of the Background Art\nGenerally, a sample observation is carried out with the use of a microscope. This sample observation is made as follows. An observation sample is placed on a microscope stage and an objective lens is moved closer to the observation sample. By doing so, an observation spot of the observation sample is observed under a magnified state.\nIn this case, an objective lens has its focal depth decreased as the magnifying power becomes higher. It is, therefore, difficult to achieve a focus setting between the objective lens and the observation sample. Further, if the distance between the objective lens and the observation sample minutely varies, defocusing occurs between the focal point position of the objective lens and the observation sample due to a variation of the minute distance, so that the quality of the observation image is greatly degraded.\nOn the other hand, an apparent position between the objective lens and the observation sample is very closer to each other. By the way, there exists a mechanical coupling length between the objective lens and the observation sample. This mechanical coupling length is constituted by many mechanical component parts present between the objective lens and the observation sample. The mechanical coupling length is provided by the length over which, for example, a microscope frame, objective lens moving mechanism, objective lens mounting revolver are passed. Therefore, the length is very long because many mechanical components are interposed.\nThese mechanical component parts are liable to be varied in their dimensions due to a temperature variation involved.\nFurther, the greater the number of the mechanical component parts the greater the mechanical coupling length involved. The microscope is easily affected by a vibration and the vibration amplitude becomes greater.\nIf, therefore, the ambient temperature varies due to, for example, the turning ON/OFF of an illumination, internal power supply, etc., as well as the operation of an air conditioning equipment, then there arises a variation in the dimensions of mechanical component parts in the microscope. Even if, therefore, the focal setting of the objective is made relative to the observation sample, the distance between the objective lens and the observation sample greatly varies due to the above-mentioned dimensional variation and there occurs a defocusing.\nFurther, under a somewhat smaller external vibration, a greater vibration amplitude is involved and a distance between the objective lens and the sample varies, thus resulting in an out-of-focus state.\nHeretofore, various kinds of autofocusing mechanisms have been considered to compensate such defocusing. These autofocusing mechanisms require a complex mechanical/electrical mechanism and control system. For this reason, the resultant apparatus becomes bulkier and expensive. Further, another microscope is known in which, like a fluorescent observation, the brightness is extremely darker at a time of observation. To such a microscope it is originally difficult to apply the above-mentioned autofocusing mechanism.\nThe following technique is disclosed in JPN PAT APPLN KOKAI Publication No. 9-120030. In this technique, a stage is provided through a rack and drive gears and it is driven in the optical axis direction of the objective lens. Between the gears and the stage, at least two rods are inserted. These rods are different in their thermal expansion coefficients and configured such that the direction of the thermal expansion coefficient of one rod acts in a direction opposite to that of the other rod. By doing so, defocusing is compensated.\nIn such a technique, the rods are arranged within the microscope body and a time is taken until the temperature is compensated under a variation of the ambient temperature. For this reason, there is a risk that, when the sample is observed, its operation, efficiency will be lowered. Further, since the rods are interposed, the mechanical coupling length becomes greater and its structure is liable to be affected from the external vibration. It is, therefore, necessary to remodel the microscope itself.\nIn the case where a living cell, etc., is observed as a target object, only a periphery side of the observation sample is sometimes warmed by a warmer. However, it is not possible to compensate a temperature drift involved.\nIt is accordingly the object of the present invention to provide a focus stabilizing apparatus which can make a focus setting of an objective lens relative to an observation sample under a stable way without being adversely influenced from a variation in the ambient temperature.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to a observation sample, and position adjusting means provided at an outer periphery of the objective lens and adapted to move the observation sample in an optical axis direction to set a focal point of the objective lens to the observation sample.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to the observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens and adapted to move the observation sample in an optical axis direction through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set through the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an objection lens, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and adapted to move the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base for placing the observation sample thereon, and an operation ring threadably inserted over the fixing base and, by being rotated relative to the fixing base, moving in the optical axis direction of the objective lens and moving the sample retaining base in the optical axis direction of the objective lens to adjust a relative distance between the observation sample and the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens to move the observation sample in the optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base placing the observation sample thereon, and an operation ring threadably inserted over the fixing base and, by being rotated relative to the fixing base, moving in the moving axis direction of the objective lens and moving the sample retaining base in the optical axis direction of the objective lens to adjust a relative distance between the observation sample and the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on a stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens, and a sample retaining base placing the observation sample thereon and threadably inserted over the fixing base and, being rotated relative to the fixing base, moving in the optical axis direction of the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens and a sample retaining base placing the observation sample thereon and threadably inserted over the fixing base and, by being rotated relative to the fixing base, moving in the optical axis direction of the objective lens.\nIn such focus stabilizing apparatus, the objective lens is fixed to the fixing base and threadably inserted directly into a revolver provided at a microscope body.\nThe position adjusting means has a rotation stop section arranged at the upper side of the sample retaining base, an intermediate seat arranged on the upper side of the rotation stop section and having the observation sample placed thereon, and a pin having one end mounted in the fixing base and the other end inserted through a hole in the rotation stop section and provided parallel to the optical axis direction of the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base placing the observation sample thereon, fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens, and a friction member situated between the fixing base and the sample retaining base and retaining the sample retaining base by a friction force relative to the fixing base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain an observation sample, position adjusting means provided on an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base placing the observation sample thereon and fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens, and a friction member situated between the fixing base and the sample retaining base and retaining the sample retaining base by a friction force relative to the fixing base.\nIn the focus retaining apparatus, a sample retaining base grasping mechanism is provided for grasping and fixing the sample retaining base.\nFurther, a temperature sensor is provided on the fixing base to detect the ambient temperature and a temperature adjusting means is provided for adjusting the temperature of the fixing base and sample retaining base on the basis of a detection output of the temperature sensor to be made constant at all times.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens, a piezoelectric actuator provided on the upper side of the sample retaining base and performing its extending/contracting operation in the optical axis direction of the objective lens, a sample base provided on the upper side of the piezoelectric actuator and placing the observation sample thereon and an electrostatic sensor provided on the upper side of the sample retaining base to detect a moving amount of the sample base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens, a piezoelectric actuator provided on the upper side of the sample retaining base and performing its extending/contracting operation in the optical axis direction of the objective lens, a sample base provided on the upper side of the piezoelectric actuator and placing the observation sample thereon, and an electrostatic sensor provided on the upper side of the sample retaining base to detect a moving amount of the sample base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens and having a flange at its outer periphery and a sample retaining base fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens and having a flange at its outer periphery, and a feed screw section provided between the flange of the fixing base and the flange of the sample retaining base to feed the sample retaining base in the optical axis direction of the objective lens relative to the fixing base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation stage, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens and having a flange at its outer periphery, a sample retaining base fitted over the outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens and having a flange at its outer periphery, and a feed screw section provided between the flange of the fixing base and the flange of the sample retaining base to feed the sample retaining base in the optical axis direction of the objective lens relative to the fixing base.\nIn the focus stabilizing apparatus, leaf springs are provided on the sample retaining base to fix the observation sample thereon.\nFurther, a mechanical coupling length between the objective lens and observation sample is set by the fixing base, sample retaining base and operation ring.\nA mechanical coupling length between the objective lens and the observation sample is set by the fixing base and the sample retaining base.\nIn the focus stabilizing apparatus, the fixing base, operation ring and sample retaining base are made of at least two different kinds of materials and selectable in their dimensions.\nThe operation ring and sample retaining base are made of at least two different kinds of materials and selectable in their dimensions.\nThe operation ring and sample retaining base are formed of materials of different linear expansion coefficients.\nAn auxiliary member is located between the operation ring and the sample retaining base and made of a material different in liner expansion coefficient from the materials of the operation ring and sample retaining base.\nThe elastic means is comprised of leaf springs provided on the stage and arranged along the direction orthogonal to the optical axis of the objective lens and a magnet provided on the leaf spring to attract the sample base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens to set a focal point of the objective lens to the observation sample, a minute movement mechanism for minutely displacing the objective lens in the optical axis direction of the objective lens, displacement amount detecting means for detecting a displacement amount of the objective lens, and control means for operating the minute movement mechanism on the basis of a detection output of the displacement amount detecting means to control a relative distance between the objective lens and the observation sample to a predetermined distance.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, a minute movement mechanism for minutely displacing the objective lens in the optical axis direction of the objective lens, displacement amount detecting means for detecting a displacement amount of the objective lens, control means for operating the minute movement mechanism on the basis of a detection output of the displacement amount detecting means to control a relative distance between the objective lens and the observation sample to a predetermined distance, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage.\nIn the focus stabilizing apparatus, the position adjusting means has a fixing base for fixing the objective lens through the minute movement mechanism, a sample retaining base for placing the observation sample thereon, and an operation ring threadably inserted over the fixing base and, by being rotated relative to the fixing base, moving in the optical axis direction of the objective lens and moving the sample retaining base in the optical axis direction of the objective lens to adjust a relative distance between the observation sample and the objective lens.\nA mechanical coupling length between the objective lens and the observation sample is set by the fixing base for fixing the objective lens through the minute movement mechanism, the sample retaining base for placing the observation sample thereon and the operation ring inserted over the fixing base and, by being rotated relative to the fixing base, moving in the optical axis direction of the objective lens and moving the sample retaining base in the optical axis direction of the objective lens to adjust a relative distance between the observation sample and the objective lens.\nThe minute movement mechanism has a moving stage relative to which the objective lens is provided, and an actuator for minutely moving the moving stage in the optical axis direction of the objective lens.\nThe minute movement mechanism has a moving stage relative to which the objective lens is provided, and piezoelectric actuators arranged in those positions symmetric relative to the optical axis of the objective lens to minutely move the moving stage in the optical axis direction of the objective lens.\nThe displacement amount detecting means is comprised of an electrostatic capacity sensor.\nThe control means has means for receiving, as inputs, an objective lens displacement amount detected by the displacement amount detecting means and an instruction value representing the position of the objective lens, means for finding a deviation between the displacement amount of the objective lens and the instruction value, and means for performing operation control of the minute movement mechanism in accordance with the deviation to move the objective lens to a position designated by the instruction value.\nThe elastic means is comprised of leaf springs mounted on the stage and arranged along the direction orthogonal to the optical axis of the objective lens and a magnet mounted on the leaf spring to attract the sample base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a minute movement mechanism for minutely displacing the objective lens in the optical axis direction of the objective lens, displacement amount detecting means for detecting a displacement amount of the objective lens, and control means for operating the minute movement mechanism on the basis of a detection output of the displacement amount detecting means to control a relative distance between the objective lens and the observation sample to a predetermined distance.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a fixing base for fixing the objective lens, a minute movement table on which the objective lens is provided, parallel springs situated between the fixing base and the minute movement table to allow the minute movement table to be moved in an optical axis direction of the objective lens, an actuator provided between the fixing base and the minute movement table to minutely displace the minute movement table in the optical axis direction of the objective lens, a displacement sensor for detecting a displacement amount of the objective lens, and control means for allowing the actuator to perform its extending/contracting operation on the basis of a detection output of the displacement sensor to control the objective lens and bring it to a just-in-focus position relative to the observation sample.\nIn the focus stabilizing apparatus, the control means has a memory section for storing an output of the displacement sensor corresponding to a just-in-focus state between the observation sample and the objective lens, a comparing section for comparing an output of the displacement sensor and an output of the displacement sensor stored in the memory section, and a control section for outputting an electric signal for canceling a distance variation between the observation sample and the objective lens on the basis of a result of comparison by the comparing section to the actuator.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens to set a focal point of the objective lens to the observation sample, a minute movement mechanism for minutely varying a distance between the observation sample and the objective lens in the optical axis direction, and control means for operating the minute movement mechanism to make a fine adjustment of the distance between the observation sample and the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises a fixing base having an objective lens mounting section, a sample base for supporting an observation sample, and position adjusting means provided on the fixing base to allow the sample base to be moved in an optical axis direction of the objective lens.\nIn the focus stabilizing apparatus, the position adjusting means has a fixing base for fixing the objective lens, an operation ring inserted over the fixing base and, by being rotated, moving in the optical axis direction of the objective lens, and a sample retaining base so provided as to be movable in the optical axis direction and, upon a rotation operation of the operating ring, adjusting a distance between the observation sample and the objective lens.\nAccording to the present invention, the mechanical coupling length between the objective lens and the observation sample is determined by only the position adjusting means provided relative to the objective lens and can be set to be very short. By doing so, there is almost no variation in a positional relation between the objective lens and the observation sample even if the ambient temperature varies. The structure suffers no adverse effect from an external vibration and the sample can be observed at all times under a stable condition.\nAfter the focal point of the objective lens has been set to the observation sample, the observation sample can be moved in a direction orthogonal to the optical axis of the objective lens simply by moving the stage and it is also possible to readily position the observation sample in a moving direction of the stage.\nFurther, as the position adjusting means use is made of at least two different kinds of materials and its component elements are selectable in their materials and dimensions. It is, therefore, possible to prevent an observation image from being degraded even if the ambient temperature varies.\nFurther, the mechanical coupling length between the objective lens and the observation sample is set to be short and the operation of the minute movement mechanism is controlled in accordance with a deviation between the displacement amount of the objective lens and the instruction value to allow the objective lens to be set to a position designated by the instruction value. Therefore, even if the ambient temperature varies, there is almost no variation in positional relation between the objective lens and the observation sample. The structure also suffers no adverse effect from external vibrations and the objective lens can be set to a desired position given by the instruction value and the observation sample can be observed, under a stable way, over an extended period of time.\nAdditional objects and advantages of the invention will be set forth in the description which follows, and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention may be realized and obtained by means of the instrumentalities and combinations particularly pointed out hereinafter."}
-{"text": "1. Field of the Invention\nThe present invention relates to an inkjet recording apparatus which ejects ink droplets to record an image on a recording medium.\n2. Description of the Related Art\nJapanese Unexamined Patent Publication No. 59597/2005 (Tokukai 2005-59597) discloses an inkjet printer including a conveyance mechanism having a drum which rotates to convey a sheet carried on the outer circumferential surface thereof; a plurality of inkjet heads each having an ejection surface, which are aligned in a conveyance direction of the sheet so that the ejection surface of the each inkjet heads faces the outer circumferential surface of the drum; and a wiper for wiping the ejection surface. In this inkjet printer, all the inkjet heads are fixed on a frame structure. The frame structure is moveable between a printing position and a wiping position. The printing position is a position where the frame structure is disposed when ink droplets are ejected from an ejection surface to a sheet conveyed by the conveyance mechanism. The wiping position is such a position that the ejection surface is disposed farther apart from the outer circumferential surface of the drum, compared to the printing position. At a time of printing, the wiper is in a standby position and faces no ejection surfaces, and the frame structure is positioned in the printing position. At a time of a wiping operation, the frame structure moves to the wiping position. Then, the wiper moves in the circumferential direction of the drum, from the standby position to an opposing position so as to face the ejection surface. Then, the wiper reciprocates in the axial direction of the drum, thereby wiping the ejection surface."}
-{"text": "This invention relates to a method and control for defrosting the outdoor coil of a heat pump in a manner which optimizes efficiency and conserves energy.\nWhen a heat pump operates in its heating mode, frost builds up on the pump's outdoor coil. As the frost thickness increases, heat transfer from the outdoor air decreases and the efficiency of the heat pump drops significantly, a substantial amount of energy therefore being wasted. Hence, it is necessary to periodically defrost the outdoor coil. This is usually accomplished by reversing the refrigerant flow in the heat pump which will heat the outdoor coil and melt the frost.\nIt is recognized that there is an optimum point of frost accumulation at which the heat pump should be switched to its defrost mode of operation. If defrost is commanded too soon or too late, energy will be wasted and efficiency will suffer. Unfortunately, it has been very difficult to achieve such optimum operation in the past. Moreover, these previous defrost systems are unreliable in operation and/or are not adaptable to all types of outdoor coils.\nSubstantially less expensive defrost control systems have also been developed, but these systems are not capable of adjusting to the prevailing weather conditions. In one such system, the differential between the outdoor ambient (dry bulb) temperature and the refrigerant temperature in the outdoor coil is measured. The outdoor coil temperature decreases as frost builds up, and this increases the temperature split or difference between the outdoor ambient temperature and the coil temperature. When the temperature split increases to a predetermined value, the outdoor coil is defrosted. These prior temperature differential type defrost controls, however, fail to take the weather conditions into account. The temperature split between the outdoor ambient air (dry bulb) temperature and the refrigerant temperature in the outdoor coil for clean coil operation is a function of the outdoor wet bulb temperature and not the dry bulb temperature. For example, when the outdoor ambient air has a 35.degree. F. dry bulb temperature, a 34.degree. F. wet bulb temperature, and a relative humidity of about 90%, the refrigerant temperature in the outdoor coil of a typical three ton heat pump may be about 23.degree. F. when the outdoor coil is frost-free, the clean coil temperature split (namely, the outdoor ambient temperature minus the outdoor coil temperature) thereby being 35.degree.-23.degree. or 12.degree.. (All temperatures mentioned herein will be F or Fahrenheit.) For the same outdoor dry bulb temperature, an outdoor wet bulb temperature of 28.degree. and an outdoor relative humidity of about 40% may then provide an outdoor coil temperature of about 17.degree., resulting in a clean coil temperature split of 35.degree.-17.degree. or 18.degree.. Neither humidity condition is uncommon in most areas. Thus, if the defrost control were set, when the ambient air has a 34.degree. wet bulb temperature, to initiate defrost at a temperature differential of, for example, 5.degree. above its expected clean coil condition, defrost would occur when the temperature differential became 12.degree.+5.degree. or 17.degree. and dry weather conditions would result in the system continually defrosting itself without time for frost buildup on the outdoor coil.\nEven if the temperature split, at which defrost should occur, is properly determined when the outdoor coil is frost-free, long before frost builds up and that temperature split is reached the weather conditions (namely, the outdoor temperature and/or relative humidity) may change significantly, and that previously determined temperature split may no longer be appropriate or valid. If there is a decrease in outdoor temperature between defrost modes, excessive frost would build up on the outdoor coil and defrost should now be initiated at a smaller temperature split, not the one previously determined. On the other hand, as the outdoor temperature rises the same system may go into needless defrost because the control would assume that frost is building up on the coil, when it may not.\nThis phenomenon may be appreciated and more fully understood by observing FIG. 1 which provides a graph of the performance of the typical three ton heat pump mentioned previously. The graph plots the wet bulb temperature of the outdoor air versus the outdoor ambient or dry bulb temperature at different outdoor relative humidities. The graph shows the liquid line temperature, which is essentially the same as the outdoor coil temperature or the coil surface temperature, under clean coil conditions at various wet bulb temperatures. The clean coil temperature splits (the outdoor dry bulb temperature minus the liquid line temperature) for different weather conditions, namely at different points on the graph, may easily be determined by subtraction of one temperature from the other at the point that represents the weather conditions. The graph clearly illustrates that the liquid line temperature is strictly a function of the wet bulb temperature, and thus the moisture in the outdoor air.\nIt will be assumed that on a given day at about 7 a.m. the weather conditions in a particular area are as depicted by point 11 in FIG. 1, namely about 12.degree. outdoor ambient temperature, 10.5.degree. wet bulb temperature and about 77% relative humidity, the liquid line temperature for clean coil conditions thus being about 4.5.degree. to provide a clean coil temperature split of 12.degree.-4.5.degree. or 7.5.degree.. Point 12 indicates the assumed weather conditions on the same day at 10 a.m.--29.degree. outdoor dry bulb temperature, 23.degree. wet bulb temperature, about 40% relative humidity and a liquid line temperature of about 13.5.degree., the clean coil temperature split thereby being 29.degree.-13.5.degree. or 15.5.degree.. This corresponds to an 8.degree. increase (15.5-7.5) in the temperature split for a clean outdoor coil. If the control system were programmed, in accordance with the data at 7 a.m., to initiate defrost after there is a 4.degree. temperature increase in the clean coil temperature split, a needless defrost cycle would occur with no frost build up on the outdoor coil. Points 13 and 14 in FIG. 1 depict the assumed weather conditions at 4 p.m. and 11 p.m., respectively, on the same given day. The graph indicates that the clean coil temperature split would change downward from about 18.degree. to 11.5.degree., or about 6.5.degree., between 4 p.m. and 11 p.m. Thus, a 4.degree. programmed differential would require that the initial 18.degree. clean coil split at 4 p.m. would have to increase to 22.degree. before defrost would occur, whereas the optimum defrost split (the difference between the outdoor temperature and the coil temperature when the defrost mode should be initiated) for the weather conditions at 11 p.m. would be 11.5.degree. plus 4.degree., or 15.5.degree.. Hence, the split would increase 6.5.degree. (from 15.5.degree. to 22.degree.) above the optimum defrost condition before defrost would be initiated and excessive frost would accumulate. The conditions assumed in explaining the FIG. 1 graph are not uncommon, since the outdoor temperature and relative humidity may experience wide variations over a 24-hour period.\nThe defrost control system of the present invention is a substantial improvement over those previously developed. The system is not only relatively inexpensive but the initiation of outdoor coil defrost is timed to occur at the optimum point regardless of changing weather conditions so that defrost only and always occurs when it is necessary, thereby increasing the efficiency of the heat pump, conserving energy and improving system reliability. Any time there is a significant change in the weather conditions, the control system of the present invention will effectively recalculate when a defrost cycle should be initiated."}
-{"text": "Many aerial vehicles (e.g., manned or unmanned vehicles such as airplanes, helicopters or other airships) are configured to operate in two or more flight modes. As one example, an aerial vehicle may be configured to engage in forward flight, or substantially horizontal flight, a mode in which the aerial vehicle travels from one point in space (e.g., a land-based point or, alternatively, a sea-based or air-based point) to another point by traveling over at least a portion of the Earth. In forward flight, the aerial vehicle may be maintained aloft by one or more net forces of lift that are typically induced by airflow passing over and below wings, consistent with a pressure gradient. As another example, an aerial vehicle may be configured to engage in vertical flight, a mode in which the aerial vehicle travels in a vertical or substantially vertical direction from one altitude to another altitude (e.g., upward or downward, from a first point on land, on sea or in the air to a second point in the air, or vice versa) substantially normal to the surface of the Earth, or hovers (e.g., maintains a substantially constant altitude), with an insubstantial change in horizontal or lateral position. In vertical flight, the aerial vehicle may be maintained aloft by one or more net forces of lift that are typically induced by rotating blades of a propeller or another source. As yet another example, an aerial vehicle may be configured to engage in both forward and vertical flight, a hybrid mode in which a position of the aerial vehicle changes in both horizontal and vertical directions.\nAn aerial vehicle that is configured to operate in multiple modes may utilize one or more propulsion systems and/or control surfaces (e.g., wings, rudders, ailerons, flaps or other components) at different times, depending on requirements of a given mission in which the aerial vehicle is to operate in each of such modes. For example, an aerial vehicle may utilize a first set of motors or rotors when operating in forward flight, and a second set of motors or rotors when operating in horizontal flight. Likewise, the aerial vehicle may utilize a first set of control surfaces when operating in horizontal flight, and a second set of control surfaces when operating in vertical flight. When motors, rotors, control surfaces or other components of an aerial vehicle are not being utilized for propulsion or control, such components merely act as dead weight to the aerial vehicle.\nThe use of imaging devices or other sensors on aerial vehicles is increasingly common. In particular, unmanned aerial vehicles, or UAVs, are frequently equipped with one or more imaging devices such as digital cameras; position sensors such as Global Positioning System, or GPS, sensors; radar sensors; or laser sensors, such as light detection and ranging, or LIDAR, sensors. Such sensors aid in the guided or autonomous operation of an aerial vehicle, and may be used to determine when the aerial vehicle has arrived at or passed over a given location, when the aerial vehicle is within range of one or more structures, features, objects or humans (or other animals), or for any other purpose. Outfitting an aerial vehicle with one or more of such sensors typically requires installing housings, turrets or other structures or features by which such sensors may be mounted to the aerial vehicle. Such structures or features add weight to the aerial vehicle, and may increase the amount or extent of drag encountered during flight, thereby exacting a substantial operational cost from the aerial vehicle for the use of such sensors in exchange for their many benefits."}
-{"text": "Various mechanisms exist for power management in a platform. However, existing power management techniques require a process or agent to have access to all components to be managed.\nExisting multi-processor platforms may have processors, devices and/or memory partitioned. One platform may then behave as two computing devices. A platform may have multiple processors with hardware partitioning so that the platform may act as multiple machines with multiple operating systems. There may be only a single motherboard on the platform.\nPower utilization of platforms is becoming more important. It is desirable to minimize or optimize the power that a platform or set of platforms utilize. As partitioning becomes more prevalent, it may become desirable to manage the power utilization of the various partitions of a system and the system overall.\nHowever, partitioning has traditionally been performed in enterprise-class machines, such as servers, and other large-size deployed machines. Power utilization has become more important as more machines are used in one location. Environmental concerns are now an issue. For instance, it may not be possible to co-locate enough compute power and maintain an appropriately cooled environment. The air conditioning unit may not be able to handle the heat dissipation of that much power utilization. With the multi-core and mini-core deployments planned for the future, many usage models require partitioning. It is becoming more common for platforms to comprise more than one processor. Further, other devices are often present in multiples within the platform.\nPortions of a system may be sequestered, for partitioning. Power management in an existing system may use an agent that can monitor the processor burden rate (busy time), control throttling, etc. The agent is typically in control of the platform, but when resources are sequestered, the agent no longer maintains control of the sequestered resources. If the agent cannot see a resource, the agent can no longer manage that resource. Thus, the traditional power management schemes are no longer viable solutions for partitioned platforms with sequestering."}
-{"text": "File systems organize and track where data is stored in memory and where free or available space exists in memory. Distributed or clustered file systems can store thousands or even millions of files. These files and the corresponding metadata are distributed across a large number of storage devices, such as disk drives.\nCounters are used in the file system to track counts for various types of information. Many file systems use simple counters that are not distributed. This approach does not scale well on clustered file systems. Other file systems implement counters that cannot be recovered in the event of corruption without a scan of the entire file system by a file system checker (fsck). Such scans can be quite time-consuming for large file systems.\nMethods and systems are needed to improve the management of counters in clustered file systems."}
-{"text": "A. Technical Field\nThe present invention relates to an ink jet recording material and a coating agent therefor, more particularly, relates to: an ink jet recording material which displays sufficient water resistance and further, excellent definition of initial images; and a coating agent therefor to give such an ink jet recording material.\nInk jet recording apparatuses are machines such as printers, facsimile machines, and copiers. For the ink jet recording apparatus, U.S. Pat. No. 5,486,854 issued Jan. 23, 1996 is hereby incorporated by reference. The ink jet recording apparatus is to make a recording by jetting ink from a recording means (i.e. a recording head which is a part of the ink jet recording apparatus) onto a recording material, wherein the recording material is, for example, a paper sheet, or a plastic sheet such as transparent PET sheet.\nThe ink jet recording material, generally, suffers from a lack of water resistance. That is, an exposure to water usually will dissolve and destroy the image (the imaged ink). To prevent this, the image must be rendered water-resistant, and further, if the ink is, for example, an organic dye, then the image must also be fixed. Examples of known methods of giving the image the water resistance to thereby fix the image include arts to fix dyes with mordants and arts involving the use of adsorptive pigments. However, operations thereof are complicated, or the optimal method is different according to the composition of the ink, so the above known methods are not commonly usable means.\nIn comparison, JP-A-035090/1998 discloses a water-resistant composition for ink jet recording sheets, comprising a polymer containing an amino group (and/or quaternary ammonium salt) and a carboxyl group (and/or acid anhydride) and a crosslinking agent containing at least two oxazoline groups, as a method of not giving the image the water resistance, but carrying out a water-resisting pretreatment for the recording material. However, as to this technique, not only is the resultant water-resistification insufficient, but also there are problems on the definition of initial images.\nAn object of the present invention is to provide: an ink jet recording material which displays sufficient water resistance and further, excellent definition of initial images; and a coating agent therefor to give such an ink jet recording material.\nTo solve the above problems, the present invention provides the following:\nA first coating agent for ink jet recording materials, according to the present invention, comprises: an aqueous polymer having a carboxyl group; and a water-soluble polymer having an oxazoline group as a crosslinking agent.\nA second coating agent for ink jet recording materials, according to the present invention, comprises a polymer and a crosslinking agent, wherein the polymer has both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group.\nAn ink jet recording material, according to the present invention, has on at least one face thereof a coated and cured layer of the above present invention first coating agent and/or a coated and cured layer of the above present invention second coating agent.\nThese and other objects and the advantages of the present invention will be more fully apparent from the following detailed disclosure.\nThe first coating agent, according to the present invention, comprises: an aqueous polymer having a carboxyl group; and a water-soluble polymer having an oxazoline group as a crosslinking agent.\nThe aqueous polymer having a carboxyl group is not especially limited if it is a polymer that has a carboxyl group and further is aqueous (that is, water-soluble or water-dispersible). To obtain the polymer having a carboxyl group, a monomer having a carboxyl group is, for example, polymerized as a raw material, or a carboxyl group is introduced into a polymer (as prepared beforehand) by denaturation. To obtain the aqueous polymer, a hydrophilic monomer (available whether it has a carboxyl group or not) is, for example, used in the ratio of preferably 10 mol % or higher, more preferably 50 mol % or higher, to the entire monomer component.\nExamples of the monomer having a carboxyl group include: unsaturated monocarboxylic acids, such as acrylic acid, methacrylic acid, and crotonic acid; unsaturated dicarboxylic acids, such as maleic acid, itaconic acid, citraconic acid, and fumaric acid; and unsaturated dicarboxylic anhydrides, such as maleic anhydride, itaconic anhydride, and citraconic anhydride.\nExamples of the method, in which a carboxyl group is introduced into a polymer (as prepared beforehand) by denaturation, include a method including the step of jumping up a polymer, having an OH group in opposite terminal portions, with pyromellitic dianhydride.\nExamples of the hydrophilic monomer include: monomers having a carboxyl group; and other monomers, such as hydroxyethyl acrylate, hydroxyethyl methacrylate, vinylpyrrolidone, dimethylaminoethyl acrylate, and chloridized-triaminoethyl methacrylate.\nExamples of the aqueous polymer having a carboxyl group, as preferably usable in the present invention, include: polyvinyl alcohols having a carboxyl group (for example, anionic KEPS series made by Dai-ichi Kogyo Seiyaku Co., Ltd., and K Polymer made by Kuraray Co., Ltd.); (meth)acrylic ester copolymers (for example, Arolon made by Nippon Shokubai Co., Ltd.); vinyl ether-maleic anhydride copolymers (Gantrez AN series made by ISP); and aqueous polymers having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group. Particularly preferable ones are the aqueous polymers having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group.\nExamples of the above aqueous polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group include a polymer having a structural unit of general formula (1) below: \nwherein: each of R1 and R2, independently of each other, denotes a\nhydrogen atom or a methyl group;\nR3 denotes a divalent organic residue;\nR4 denotes a monovalent organic residue;\nZ denotes a halogen atom;\na denotes an integer of 1xcx9c1,000;\nb denotes an integer of 3xcx9c110,000; and\nn denotes an integer of 3xcx9c5,000.\nThis polymer can be synthesized by carrying out cationic polymerization of an oxazoline compound in the presence of an unsaturated halide to synthesize a polyoxazoline macromonomer having a radical-polymerizable double bond at a polymerization-initiating terminal, and then copolymerizing this polyoxazoline macromonomer and a monomer having a carboxyl group, as is illustrated by the following chemical reaction formula: \nwherein R1, R2, R3, R4, Z, a, b, and n are the same as those in general formula (1).\nExamples of the oxazoline compound, as usable in the above cationic polymerization, include 2-methyl-2-oxazoline, 2-ethyl-2-oxazoline, 2-(n-propyl)-2-oxazoline, 2-(i-propyl)-2-oxazoline, 2-(n-butyl)-2-oxazoline, 2-(i-butyl)-2-oxazoline, and 2-(t-butyl)-2-oxazoline. Preferable ones among them are compounds with not more than 3 carbon atoms in R4, of which specific examples include 2-ethyl-2-oxazoline.\nExamples of the above unsaturated halide include chloromethylstyrene, allyl chloride, epichlorohydrin, and chloroethyl vinyl ether. A preferable one among them is chloromethylstyrene.\nAs to the monomer having a carboxyl group, compounds in which R2 is a hydrogen atom or a methyl group are preferable among the above-exemplified monomers having a carboxyl group. Specific examples of such compounds include acrylic acid and methacrylic acid.\nTherefore, preferable examples of the aqueous polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group include poly(2-methyl-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-ethyl-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-(n-propyl)-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-(i-propyl)-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-(n-butyl)-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-(i-butyl)-2-oxazoline)/(meth)acrylic acid copolymers, and poly(2-(t-butyl)-2-oxazoline)/(meth)acrylic acid copolymers.\nThe above aqueous polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group is available whether it contains a structural unit other than the structural unit of general formula (1) above or not.\nIn addition, in the above general formula (1), a is an integer of 1xcx9c1,000, preferably 1xcx9c500, and b is an integer of 3xcx9c10,000, preferably 3xcx9c5,000, and n is an integer of 3xcx9c5,000, preferably 10xcx9c5,000, more preferably 2xcx9c100, still more preferably 5xcx9c500.\nThe weight-average molecular weight of the above aqueous polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group is preferably in the range of 50,000xcx9c1,000,000, more preferably 100,000xcx9c500,000.\nExamples of the water-soluble polymer having an oxazoline group as a crosslinking agent, as contained in the first coating agent according to the present invention, include a polymer containing an oxazoline group as obtained by polymerizing a monomer component which comprises an addition-polymerizable oxazoline and, if necessary, further comprises a monomer copolymerizable therewith. Examples of the addition-polymerizable oxazoline include 2-vinyl-2-oxazoline, 2-vinyl-4-methyl-2-oxazoline, 2-vinyl-5-methyl-2-oxazoline, 2-isopropenyl-2-oxazoline, 2-isopropenyl-4-methyl-2-oxazoline, 2-isopropenyl-4-ethyl-2-oxazoline, 2-isopropenyl-5-methyl-2-oxazoline, 2-isopropenyl-5-ethyl-2-oxazoline, and 2-isopropenyl-4,5-dimethyl-2-oxazoline. These may be used either alone respectively or in combinations with each other. A preferable one among them is 2-isopropenyl-2-oxazoline, because it is industrially easily available.\nThe amount of the addition-polymerizable oxazoline, as used, is not especially limited, but is preferably 5 weight % or larger, more preferably in the range of 30xcx9c60 weight %, of the entire monomer component. In the case where the amount is smaller than 5 weight %, the extent of the curing is insufficient. In the case where the amount exceeds 60 weight %, it will have a bad effect on the water resistance.\nExamples of the monomer copolymerizable with the addition-polymerizable oxazoline include: (meth)acrylic esters, such as methyl (meth)acrylate, butyl (meth)acrylate, 2-ethylhexyl (meth)acrylate, methoxypolyethylene glycol (meth)acrylate, polyethylene glycol mono(meth)acrylate, 2-hydroxyethyl (meth)acrylate, and 2-aminoethyl (meth)acrylate and its salts; unsaturated nitriles, such as (meth)acrylonitrile; unsaturated amides, such as (meth)acrylamide, N-ethylol(meth)acrylamide, and N-(2-hydroxyethyl)(meth)acrylamide; vinyl esters, such as vinyl acetate and vinyl propionate; vinyl ethers, such as methyl vinyl ether and ethyl vinyl ether; xcex1-olefins, such as ethylene and propylene; halogen-containing xcex1,xcex2-unsaturated monomers, such as vinyl chloride, vinylidene chloride, and vinyl fluoride; and xcex1,xcex2-unsaturated aromatic monomers, such as styrene, xcex1-methylstyrene, and sodium styrenesulfonate. These may be used either alone respectively or in combinations with each other.\nTo obtain the water-soluble polymer, the ratio of the hydrophilic monomer to the entire monomer component to be polymerized is preferably 50 weight % or higher, particularly preferably in the range of 60xcx9c90 weight %. Examples of the hydrophilic monomer include addition-polymerizable oxazolines, 2-hydroxyethyl (meth)acrylate, methoxypolyethylene glycol (meth)acrylate, polyethylene glycol mono(meth)acrylate, and 2-aminoethyl (meth)acrylate and its salts, and further, sodium (meth)acrylate, ammonium (meth)acrylate, (meth)acrylonitrile, (meth)acrylamide, N-methylol(meth)acrylamide, N-(2-hydroxyethyl)(meth)acrylamide, and sodium styrenesulfonate, as are selected from among the aforementioned monomer components.\nThe aqueous polymer having a carboxyl group is contained in the first coating agent in the ratio of preferably 5xcx9c95 weight %, more preferably 10xcx9c90 weight %, still more preferably 50xcx9c90 weight %, in terms of solid content, to the entire weight of the first coating agent. The water-soluble polymer having an oxazoline group, as used as the crosslinking agent, is contained in the first coating agent in the ratio of preferably 5xcx9c95 weight %, more preferably 10xcx9c90 weight %, still more preferably 10xcx9c50 weight %, in terms of solid content, to the entire weight of the first coating agent.\nThe first coating agent may further comprise components other than the aqueous polymer having a carboxyl group and the water-soluble polymer having an oxazoline group as the crosslinking agent, if necessary. Examples of such other components include: curing catalysts, such as paratoluenesulfonic acid (PTSA); organic or inorganic fine particles; dye mordants; pigments; dispersants; and ultraviolet absorbing agents.\nThe second coating agent, according to the present invention, comprises a polymer and a crosslinking agent, wherein the polymer has both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group. Incidentally, this xe2x80x9cpolymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl groupxe2x80x9d is mentioned in detail above.\nExamples of the crosslinking agent, as used for the second coating agent, include the water-soluble polymer having an oxazoline group as mentioned in detail above, and further, melamine, aziridine, isocyanate, and epoxy.\nThe polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group is contained in the second coating agent in the ratio of preferably 5xcx9c95 weight %, more preferably 10xcx9c90 weight %, still more preferably 50xcx9c90 weight %, in terms of solid content, to the entire weight of the second coating agent. The crosslinking agent is contained in the second coating agent in the ratio of preferably 5xcx9c95 weight %, more preferably 10xcx9c90 weight %, still more preferably 10xcx9c50 weight %, in terms of solid content, to the entire weight of the second coating agent.\nThe second coating agent may further comprise components other than the polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group and the crosslinking agent, if necessary. Examples of such other components include: curing catalysts, such as paratoluenesulfonic acid (PTSA); organic or inorganic fine particles; dye mordants; pigments; dispersants; and ultraviolet absorbing agents.\nThe ink jet recording material, according to the present invention, has on at least one face thereof a coated and cured layer of the present invention first coating agent and/or a coated and cured layer of the present invention second coating agent.\nThe ink jet recording material is a sheet-shaped material as used to record images thereon with ink jet recording apparatuses. Examples of such ink jet recording materials include: paper, synthetic paper such as Tyvek (made by E.I. Du Pont DE NEMOURS and Co., Ltd.); cloths, such as canvas, clothing fabrics and non-woven composites; films or sheets of plastics such as polyvinyl chloride, polypropylene, and polyethylene terephthalate (PET).\nThe amount of the present invention first or second coating agent, as coated, is preferably in the range of 3xcx9c50 g, more preferably 5xcx9c40 g, per square meter. In the case where the amount is smaller than 3 g, water cannot sufficiently be absorbed from ink. In the case where the amount exceeds 50 g, much time and energy are necessary for drying the sheet. In addition, the coating thickness is preferably in the range of 1xcx9c50 xcexcm, more preferably 5xcx9c40 xcexcm. In the case where the coating thickness is less than 1 xcexcm, water cannot sufficiently be absorbed from ink. In the case where the coating thickness exceeds 50 xcexcm, the improvement of the ink absorbency cannot be expected very much, so there are economical disadvantages.\nThe curing temperature is preferably in the range of 50xcx9c200xc2x0 C., more preferably 80xcx9c150xc2x0 C. The curing time depends on the curing temperature, but is preferably in the range of 1xcx9c60 minutes, more preferably 1xcx9c30 minutes.\n(Effects and Advantages of the Invention):\nCoating an ink jet recording material with the present invention coating agent for ink jet recording materials can give an ink jet recording material which displays sufficient water resistance and further, excellent definition of initial images.\nHereinafter, the present invention is more specifically illustrated by the following examples of some preferred embodiments in comparison with comparative examples not according to the invention. However, the invention is not limited to the below-mentioned examples. In addition, in the examples, unless otherwise noted, the units xe2x80x9c%xe2x80x9d and xe2x80x9cpart(s)xe2x80x9d denote those by weight.\n less than Synthesis of FX-AA as an Aqueous Polymer having a Carboxyl Group and further a Structural Unit of General Formula (1) greater than \nA mixture of 297 g of 2-ethyl-2-oxazoline, 9.2 g of chloromethylstyrene (mixture of m- and p-isomers) and 252.4 g of ethanol was charged into a 1-liter autoclave, and then heated at 130xc2x0 C. for 4 hours, thus obtaining an ethanol solution of poly(2-ethyl-2-oxazoline) macromonomer with a styrene functional group at a polymerization-initiating terminal. Then, 43.2 g of acrylic acid, 3.5 g of 2,2xe2x80x2-azobis(isobutyronitrile) and 50 g of ethanol were added to this ethanol solution, and the resultant mixture was heated within the range of 105 to 135xc2x0 C. for 6 hours. The reaction mixture was cooled, thus obtaining an ethanol solution of a poly(2-ethyl-2-oxazoline) macromonomer/acrylic acid copolymer (655.3 g of FX-AA; solid content=59.6 weight %).\n less than Synthesis of Water-soluble Polymer (A) having an Oxazoline Group as a Crosslinking Agent greater than \nA mixture of 29 parts of deionized water and 1 part of V-50 (polymerization initiator made by Wako Pure Chemical Industries, Ltd.: 2,2xe2x80x2-azobis(2-amidinopropane) dihydrochloride) were charged into a flask having a stirrer, a reflux condenser, a nitrogen-introducing tube, a dropping funnel, and a thermometer, and then heated to 60xc2x0 C. under a slow nitrogen gas current. Thereto, a monomer mixture was added from the dropping funnel over a period of 1 hour, with this monomer mixture having been prepared beforehand and comprising 0.4 parts of ethyl acrylate, 5.6 parts of methyl methacrylate, 4 parts of methoxypolyethylene glycol methacrylate (NK Ester AM-90G made by Shin-Nakamura Chemical Industrial Co., Ltd.), and 10 parts of 2-isopropenyl-2-oxazoline. During the reaction, nitrogen gas was continuously run, and the temperature in the flask was kept at 60xc2x11xc2x0 C. After the end of the addition, the same temperature was maintained for 9 hours, and then the mixture was cooled, thus obtaining an aqueous solution of a polymer (water-soluble polymer (A)) containing a 2-oxazoline group, of which the nonvolatile content was 40%, and the pH was 8.3."}
-{"text": "This invention relates to fluid flow rate measurements. More specifically, this invention relates to determining fluid flow rates using measured motor speed and torque parameters.\nMeasurement of fluid flow through pipes and pumps is well known in the art. One method of determining fluid flow rates is to install gears, vanes, paddle wheel, turbines, etc., in the flow channel and determine fluid flow rate by the speed at which these devices turn. A second method is by measuring the differential pressure across a dedicated flow obstruction, such as a venturi, orifice plate, annubar, pitot tube, etc., and applying the well-known Bernoulli\"\"s principle to obtain a velocity and, consequentially, a fluid flow rate. As an example of this principle, U.S. Pat. No. 5,129,264, entitled xe2x80x9cCentrifugal Pump with Flow Measurement,xe2x80x9d issued Jul. 14, 1992, to Lorenc, and assigned to the same assignee herein discloses using differential pressure and pump speed to measure fluid flow rates. Still other methods of fluid flow rate measurement employ electrical/magnetic or sonic measurement means. For example, Mag Meters determine fluid flow rates by measuring the change in a magnetic field caused by the velocity of the fluid flowing therethrough. Sonic devices use acoustical pulses, i.e., Sonar, and Doppler principles to measure fluid flow rates. Other non-intrusive methods measure the torque a variable speed electrical motor delivers to a pump by installing a torque shaft between the motor and pump. The motor or pump is then calibrated and a motor kilowatt input/Motor Brake Horsepower Output (BHP) calibration table is developed. Accordingly, knowledge of the kilowatt input value can be used to determine the output horsepower. However, this calibration is needed at several speeds and requires several different sized torque shafts.\nThus, current methods for determining fluid flow rates necessitate intrusion into, or require access to, the enclosures transporting the fluid. In some systems, such as caustic systems having lined pumps, intrusion is prohibited.\nHence, there is a need for a simple, accurate and reliable method for determining fluid flow rates without intruding into the fluid flow or having access to the enclosures transporting the fluid.\nThe present invention determines a fluid flow rate through a pump by first determining two flow rate values from a plurality of characterizing flow rate values corresponding to two known speed values selected from a plurality of known characterizing speed values at a known motor torque. The first one of the two known speed values is selected greater than a measured pump speed and a second one of the two known speed values is selected less than the measured pump speed. The fluid flow rate is then determined as being proportional to the two determined flow values at the known speed values and the pump speed."}
-{"text": "Not applicable.\nThe present invention relates to a marine or other recreational vehicle receptacle for receiving and storing a marine shower. The marine shower may or may not accommodate a hot-cold water mixing valve as required by the consumer.\nWhile the invention has multiple uses readily apparent to those skilled in the art it will be described primarily for use in connection with the marine industry.\nMarine showers are most often mounted in the transom of a boat or in the cockpit coaming of a boat and in modern day boat designs these transoms and coamings comprise a multitude of surfaces with orientations from vertical to horizontal and therebetween. These orientations require different shower head receptacle designs and when a mixing valve is required a still further design is needed. Thus multiple items must be kept in inventory in order to accommodate all requirements.\nA shower receptacle designed for a vertical surface will not be useful for a horizontal surface and a receptacle adapted for a straight shower handle will not accommodate a Euro-style handle with a projecting on-off lever. When a hot-cold water mixing valve is desired a further receptacle is needed to accommodate the valve. Some prior art designs feature the shower wand and hot and cold water valves in one box but this requires a large opening to be cut in the boat transom with its attendant disadvantages.\nThe disadvantages and limitations of the prior art devices are obviated by the present invention.\nAn object of the present invention is to provide a shower head receptacle adapted to be mounted on surfaces having different angles of orientation.\nA further object of the present invention is to provide a shower head housing or receptacle for receiving either a straight or Euro-style shower wand or handle.\nA still further object of the present invention is to provide a receptacle with a lid which is relatively easy to manufacture.\nAnother object of the present invention is to provide a receptacle which can accommodate various shower heads or wands as well as a hot and cold water mixing valve as the case may be.\nAnother object of the invention is to provide a receptacle for a shower head and a receptacle for a mixing valve in side by side relation with a single cover for both receptacles.\nOther objects, advantages and features of the present invention will become apparent to those skilled in the art from the following detailed description which with referencess to the accompanying drawings discloses a preferred embodiment of the invention."}
-{"text": "This invention relates to variable area exhaust nozzles for gas turbine engines and, more particularly, to sealing means for nozzle flaps of turbojet engines.\nThe exhaust nozzle of a gas turbine engine, such as a turbojet or turbofan engine, has as its purpose a transformation of the pressure and thermal energy of the combustion discharge into velocity, with the forward thrust of the engine being directly proportional to the increase in velocity of the gas from the entrance of the engine to the exit plane of the nozzle. In high performance engines and, in particular, in engines having some sort of thrust augmentation such as an afterburner, it has been found desirable to cause a variation of nozzle flow area to maintain high engine performance under a wide range of operating conditions. For example, it is desirable to provide a larger nozzle flow area during a take-off mode of operation than during a cruise mode. In addition to the function of maintaining the exhaust gas temperature within allowable limits, variable area exhaust nozzles may be employed to bring about noise, thrust and fuel economy benefits. One means for varying the nozzle flow area is by the so-called iris mechanism wherein a plurality of concentrically disposed movable members or flaps are pivotably supported about the nozzle axis. One of the problems associated with such an arrangement is the need to maintain effective sealing between the flaps as the flaps are adjusted to vary the nozzle flow area. Therefore, it is desirable to provide an exhaust nozzle whose area can be flexibly varied between minimum and maximum positions while maintaining a circumferentially continuous aerodynamic structure throughout the entire range.\nEarly method of locating seals with respect to exhaust nozzle flaps relied entirely on a combination of bolts and spectacles wherein, when the nozzle was in the closed position the seals were relatively free to move in the circumferential direction, and when the flaps moved toward the open position, the position of the seals was still not positively enough controlled so as to maintain circumferential sealing integrity around the entire nozzle periphery. Some of the problems encountered were those of dimensional stack-up, limited seal overlap within the circumferential envelope, and misalignment due to nozzle sag on or near the horizontal plane. These problems caused nozzle leakage and seal \"blow-out\", thereby resulting in decreased nozzle efficiency.\nRecent methods of effecting positive placement of seals within exhaust nozzles employ a combination of linkage pairs interconnecting the flaps to the seal wherein an axial track is located on the seal for the purpose of varying the effective lengths of the links. Such an arrangement has been recognized as being somewhat complex and requiring an excess number of moving parts which are susceptible to wear and malfunction.\nAccordingly, a primary object of the present invention is to provide an improved seal arrangement for a jet engine variable exhaust nozzle flap.\nAnother object of this invention is the provision in a variable exhaust nozzle for the maintaining of circumferential sealing integrity throughout the range of nozzle areas.\nYet another object of the present invention is the provision for maintaining a variable area exhaust nozzle seal in a centered relationship between adjacent flaps during all modes of nozzle operation.\nThese objects and other features and advantages become more readily apparent upon reference to the following description when taken in conjunction with the appended drawings."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to an injection apparatus; particularly an injection apparatus maintaining the nozzle and the injection gate at respective desired temperatures.\n2. Background of the Invention\nIt has long been known that the temperature of a melt material is important to successful injection. This is particularly true when the melt material has a high melt temperature. For example, polyethylene terephthalate (\"PET\") is typically injected above 500.degree. F. A drop in the temperature of the melt material prior to reaching the injection cavity would lower the melt material temperature below that required for proper melt material flow causing less than ideal flow characteristics. These flow characteristics can cause deformed or defectively molded parts; particularly when injecting multilayer parts comprising very thin layers. Therefore, it is desirable to maintain the nozzle temperature at or above the temperature required to assure proper melt material flow as the melt material leaves the nozzle.\nIt is also known to maintain an injection cavity at a temperature relatively low compared to the temperature of the melt material to facilitate quick cooling of the melt material upon reaching the cavity. The colder the cavity temperature at the time the melt material is injected, the faster the melt material will solidify and allow removal of the solidified part from the cavity. Therefore, a relatively lower cavity temperature will decrease the overall cycle time for injection molding a part. Moreover, it is known that if the injection gate temperature exceeds the desired temperature of the melt material, `stringing` of the melt material will occur in the nozzle and gate area as the injected part is removed from the cavity after injection is complete. These `strings` either break off with the injected part and interfere with further processing of the part (e.g. blowmolding) or stay in the gate or cavity and cause a physical or aesthetic defect in subsequently injected parts.\nFor these reasons, it has been found desirable to prevent excessive heat transfer from the injection nozzle to the injection cavity. The melt material can thus be maintained at its appropriate temperature in both the nozzle and the cavity. Prior injection apparatuses were often designed to space a nozzle tip from an associated injection cavity during injection to leave a gap therebetween. It was thought that this gap would act as a thermal break between the nozzle and the cavity and allow the nozzle to operate at high temperatures while maintaining a relatively cool cavity. Unfortunately, the thermal break of this configuration could not be maintained at efficient cycle times. During the injection process, melt material would deviate from the injection path and flow into the gap between the nozzle and the cavity. The thermal break thus became a thermal bridge.\nOther attempts to insulate an injection nozzle from a cavity have involved the use of nozzle inserts. For example, U.S. Pat. No. 4,279,588 issued to Gellert and entitled \"Hot Tip Seal\" disclosed a seal (12) located between the nozzle and the injection gate to limit heat transfer therebetween. The seal (12) of Gellert resided substantially within the nozzle and extended outward therefrom to contact the cavity. Similarly, U.S. Pat. No. 4,521,179 issued to Gellert and entitled \"Injection Molding Core Ring Gate System\" disclosed a nozzle seal (76). The seal (76) of Gellert also resided substantially within the nozzle and extended outward therefrom to contact the cavity.\nIt has been found that movement of the various parts within an injection apparatus will result from thermal expansion as portions of the apparatus are heated from ambient temperature to the temperature necessary to inject a melt material. Different injection apparatuses accommodate this thermal expansion in different ways. It has been found that the thermal expansion of some injection apparatuses results in movement of the nozzle both along the longitudinal axis thereof and perpendicular to that longitudinal axis. In other words, it has been found that the nozzles of some apparatuses will elongate and shift laterally as the apparatus is heated. Seals that attached to the nozzle, such as those of the Gellert patents discussed above, break or deform due to this lateral nozzle movement. Such seals are therefore inapplicable to apparatuses experiencing this lateral nozzle movement.\nIt has also been found that many seals cannot withstand the high temperatures and pressures associated with injection; especially when the high temperatures are maintained for long periods of time. Many prior inserts degraded after prolonged exposure to high temperatures resulting in rupture or deformation of the inserts which allowed melt material to leak into the area between the nozzle and the cavity causing in a thermal bridge.\nIt has also been known to supply a cooling means to a cavity to remove the heat transferred from the nozzle or melt material to the cavity. Cooling ducts circulating coolants such as glycol were typically employed. However, the distance between the part void and the injection gate has heretofore limited the proximity of the cooling ducts to the injection gate."}
-{"text": "Not applicable\nThe present invention is directed to an OFDM/DMT digital communications system. More particularly, the present invention is directed to an apparatus and method for synchronizing the clocks used in a transmitter and receiver of an OFDM/DMT digital communications system. The present invention is particularly applicable in multipoint OFDM/DMT digital communications systems.\nMulti-point communications systems having a primary site that is coupled for communication with a plurality of secondary sites are known. One such communications system type is a cable telephony system. Cable telephony systems transmit and receive telephone call communications over the same cable transmission media as used to receive cable television signals and other cable services.\nOne cable telephony system currently deployed and in commercial use is the Cablespan 2300 system available from Tellabs, Inc. The Cablespan 2300 system uses a head end unit that includes a primary transmitter and primary receiver disposed at a primary site. The head end unit transmits and receives telephony data to and from a plurality of remote service units that are located at respective secondary sites. This communication scheme uses TDM QPSK modulation for the data communications and can accommodate approximately thirty phone calls within the 1.9 MHz bandwidth typically allocated for such communications.\nAs the number of cable telephony subscribers increases over time, the increased use will strain the limited bandwidth allocated to the cable telephony system. Generally stated, there are two potential solutions to this bandwidth allocation problem that may be used separately or in conjunction with one another. First, the bandwidth allocated to cable telephony communications may be increased. Second, the available bandwidth may be used more efficiently. It is often impractical to increase the bandwidth allocated to the cable telephony system given the competition between services for the total bandwidth available to the cable service provider. Therefore, it is preferable to use the allocated bandwidth in a more efficient manner. One way in which the assigned bandwidth may be used more efficiently is to use a modulation scheme that is capable of transmitting more information within a given bandwidth than the TDM QPSK modulation scheme presently employed.\nThe present inventors have recognized that OFDM/DMT modulation schemes may provide such an increase in transmitted information for a given bandwidth. Such systems, however, present a number of technical problems. One such problem is the determination of how one or more remote receivers are to synchronize their internal clocks and timing systems with the internal clock and timing system of a primary transmitter at a central site. A remote receiver must first synchronize its internal clock and timing system with the clock used by the primary transmitter to synthesize the transmitted signal before the remote receiver can properly demodulate the data that it receives. A further problem occurs in multipoint communication systems in which there are plural groups of remote transmitters that transmit data to a central transceiver. Each group of transmitters often has its transmissions frequency multiplexed with transmissions from other groups before being demultiplexed for receipt by the central transceiver. The resulting multiplexing/demultiplexing operations introduce frequency offsets for which compensation must be made if the receiver of the central transceiver is to properly extract the correct data from the signals that is receives. The present inventors have recognized the need for such upstream and downstream clock synchronization and have disclosed solutions to these problems.\nIn a communications system comprising a transmission medium, symbols are generated in a predetermined number of bins using at least a first timing signal. A reference signal also is generated by using said first timing signal. The symbols and said reference signal are transmitted across a transmission medium carried by a carrier signal. The generating the transmitting are preferably achieved with a transmitter. At a first point along the transmission medium, the carrier signal with a second signal are frequency multiplexed, preferably by a frequency multiplexer. The carrier signal and the second signal are transmitted to a second point on the transmission medium displaced from the first point. The carrier signal is frequency demultiplexed after the carrier signal and second signal have reached the second point, preferably by a demultiplexer. The demultiplexed carrier signal is frequency demodulated in response to the reference signal, preferably by a first demodulator. The symbols and the reference signal are demodulated in response to the frequency demodulated carrier signal, preferably by a second demodulator.\nOther features and advantages of the present invention will become apparent upon review of the following detailed description and accompanying drawings."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to integrated circuit semiconductor memory devices and associated methods. In particular, the present invention relates to a column select line (CSL) control for which the same signal controls the enable timing and the disable timing signals for synchronous random access memory devices.\n2. Description of Related Art\nSpeed improvements in semiconductor memory devices, such as Dynamic RAMs and Static RAMs, have historically come from process and photolithography advances. More recent memory speed improvements, however, have resulted mainly from making changes to the base architecture. An example of fast RAM architecture is the synchronous architecture. One key advancement of the synchronous memories is their ability to synchronously burst data at a high-speed data rate. Additionally, in a system with a synchronous RAM, since data, addresses and control signals are latched into the memory in synchronism with the system's clock signal, the system's processor is able to perform other tasks freely until data is available after a known number of clock cycles. This architecture provides substantial advantages in memory operating performance.\nIn a typical semiconductor memory device, in order to write/read data into/from a specific memory cell in a memory device, the specific memory cell should be designated by a row address and a column address. When the specific memory cell is designated in a read/write operation, a charge distribution operation is performed with respect to data read out from the designated memory cell to a bit line, and the readout data is amplified by a sense amplifier. The amplified data is transmitted to an input/output line through an I/O gate circuit, and then is output from the memory chip via associated output circuits. The read operation of one-bit data stored in the specific memory cell is completed by the above-described process. The column decoder turns on the selected I/O gate by receiving and decoding the column address.\nTo simplify the complexity of the decoding operation in highly integrated memories, a column pre-decoder is typically provided to pre-decode the column address prior to the main decoding operation therefor. This column decoding scheme has been adopted in most high density memory devices.\nFIG. 1 is a block diagram illustrating a conventional exemplary synchronous memory device. Referring to FIG. 1, an array 100 of memory cells is provided to store data. Word lines WL0-WLm and bit lines BL0-BLn coupled with the cells run along the rows and columns of the memory cell array 100, respectively. In the vicinity of the cell array 100, a row decoder 120 is provided for selectively driving the word lines WL0-WLm, and an input/output (I/O) gate circuit 140 for supporting the selective transmission of data from the bit lines BL0-BLn to a data I/O buffer 280, and vice versa. The I/O gate circuit 140 is controlled by column select lines CSL0-CSLn. Externally applied address signals A0-Ax including both column and row address signals are fed to an address buffer 160. The column address signals CA0-CAi among the address inputs are applied to a column pre-decoder 180.\nA clock buffer 230 is suppled with an external clock signal XCLK and provides an internal PCLK synchronized with the external clock signal XCLK. A CSL enable control circuit 240 generates a CSL enable control clock signal PCSLE by the logical combination of the internal clock signal PCLK and a column address setting signal PYE from a timing control logic (not seen). The column pre-decoder 180 pre-decodes the column address signals CA0-CAi and generates pre-decoded address signals DCA0-DCAj.\nThe column pre-decoder 180 outputs the DCA0-DCAj signals under the control of the PCSLE signal from the CSL enable control circuit 240. Main decoding operation of the column address signals are then carried out by a column main-decoder 200. This decoder 200 generates decoded signals DCAB0-DCABk by decoding the DCA0-DCAj. The DCAB0-DCABk signals are provided to a column driver 220 which drives the column select lines CSL0-CSLn selectively in response to the DCAB0-DCABk signals. A CSL disable control circuit 260 generates a CSL disable control clock signal PCSLD by the logical combination of the internal clock signal PCLK and a normally logic-high signal PVCCH. The column driver 220 is disabled by the PCSLD signal from the CSL disable control circuit 260, and hence stops driving the column select lines CSL0-CSLn.\nFIGS. 2A and 2B illustrate the constructions of the CSL enable and disable control circuits 240 and 260, respectively, in detail. Referring first to FIG. 2A, the CSL enable control circuit 240 includes a delay circuit formed by inverters IV1-IV4 (\"first\" delay circuit), a NAND gate G1, and an inverter IV5. The internal clock signal PCLK is provided to the delay circuit. The NAND gate G1 has one input applied with the delayed signal of the clock signal PCLK and the other input applied with the column address setting signal PYE. The output signal of the NAND gate G1 is output through the inverter IV5 as the CSL enable control clock signal PCSLE.\nReferring to FIG. 2B, the CSL disable control circuit 260 includes a delay circuit formed by inverters IV6-IV8 (\"second\" delay circuit) and a NAND gate G2. The second delay circuit is also fed with the clock signal PCLK. This delay circuit has a smaller delay time than the first delay circuit. The output of the second delay circuit is supplied to one input of the NAND gate G2. The normally logic-high signals PVCCH is provide to the other input of the NAND gate G2. This gate G2 outputs the CSL disable control clock signal PCSLD.\nFIG. 3 shows the detailed configuration of a unit circuit of the column pre-decoder 180. As shown in FIG. 3, the unit pre-decoder circuit 180' includes inverters IV31-IV49 and NAND gates G34-G49. The unit column pre-decoder circuit 180' is provided with three column address signals CA0-CA2 from the address buffer 160, and generates eight pre-decoded column address signals DCA0-DCA7. The CSL enable control clock signal PCSLE is commonly applied to the first inputs of the NAND gates G42-G49. The second inputs of the NAND gates G42-G49 are provided with the substantial pre-decoded column address signals, i.e., the output signals of the inverters IV34-IV41, respectively. When the PCSLE signal becomes high, the output signals of the inverter IV34-IV41 can be propagated to the inverters IV42-IV49 via the NAND gates G42-G49, respectively, and they are output as the pre-decoded column address signals DCA0-DCA7. The PCSLE signal should go high only after the completion of the pre-decoding operation with the inverters IV31-IV41 and the NAND gates G34-G41 in order to prevent pre-decoding errors.\nFIG. 4 illustrates the detailed construction of unit circuits of the column main-decoder 200 and column driver 220, respectively. Referring to FIG. 4, the unit column main-decoder circuit 200' includes NAND gates G50-G57, and inverters IV50-IV57 corresponding to the NAND gates G50-G57, respectively. Each of the NAND gates G50-G57 has one input provided with a corresponding pre-decoded column address signal DCAy (where, y=0, 1, . . . , or 7) and the other input with a gate control signal GCS from a timing control logic (not shown). Each output signal of the NAND gates G50-G57 is provided as a finally decoded signal DCABy (where, y=0, 1, . . . , or 7) through the corresponding inverter IV50, IV51, . . . , or IV57.\nThe unit column driver circuit 220' includes inverters IV60-IV67, cascode inverters (sometimes called \"dual gate inverters\") 40-47, and inverting latches 60-67. Each of the cascode inverters 40-47 consists of two PMOS transistors (e.g., MP40a and MP40b) and one NMOS transistor (e.g., MN40), and each inverting latch (e.g., 60) is formed of two cross-coupled inverters (IV60a and IV60b). For each cascode inverter (e.g., 40), three transistors (MP40a, MP40b and MN40) have their source-drain paths coupled in series between a boosted supply voltage terminal VEXT and a ground voltage terminal GND. Each decoded column address signal (e.g., DCAB0) is applied to the gates of the corresponding pull-up and pull-down transistors (MP40a and MN40). The PCSLD signal from the CSL disable control circuit 260 is commonly fed to the gates of the switching transistors MP40b-MP47b of the respective cascode inverters 40-47 via the inverters IV60-IV67. Each inverting latch (e.g., 60) is coupled to the drain junction of the corresponding switching and pull-down transistors (MP60a and MN60).\nThe pre-decoding operation begins with the CSL disable control clock signal PCSLD of a high level. A pre-decoded column address signal of a high level (e.g., DCA0) can be transferred to the gate of the PMOS transistor MP40a as a decoded column address signal DCABO only when the gate control signal GCS remains at a high level. Namely, the GCS signal determines whether to propagate the DCA0-DCA7 signals through the unit column main-decoder circuit 200' or not. The decoded signal DCAB0 goes high when both DCA0 and GCS signals are high, so PMOS pull-up transistor MP40a turns off and NMOS pull-down transistor MN40 on. The high-level DCAB signal is latched by the inverting latch 60, so that a corresponding column select line CSL0 is driven high. After the PCSLD signal has gone low, the GCS signal also goes low. Accordingly, the pull-up transistor MP40a turns on and the pull-down transistor MN40 off, but the column select line CSL0 still remains high owing to the inverting latch 60. In this situation, when the PCSLD signal goes high again, the switching transistor MP40b turns on, so the CSL0 line is driven low.\nAs described above, the column select lines CSL0-CSLn are selectively activated by the column pre-decoder 180, but deactivated by separately controlling the column driver 220.\nFIG. 5 is a timing diagram illustrating read/write operations of the conventional synchronous memory device of FIG. 1. With reference to FIG. 5, after a column address strobe signal CAS is activated low, in clock cycle T0, the CSL disable control clock signal PCSLD goes high in synchronism with the external clock signal XCLK (or the internal clock signal PCLK). After a predetermined time (i.e., Tm1) has elapsed, during which the first column address signals CA#0 (i.e., CA0-CAi) have reached the column pre-decoder 180, the CSL enable control clock signal PCSLE goes high in response to the activation of the column address setting signal PYE (see FIG. 2A). Of course, the PCSLE signal is also synchronized with the clock signal XCLK (or PCLK). A unit column pre-decoder circuit 180' pre-decodes the column address signals CA#0 (CA0-CA2) and generates the pre-decoded column address signals DCA#0 (DCA0-DCA7) of which only one is active and the others inactive. Here, assuming DCA0 signal is activated high, then a corresponding column select line CSL0 will be driven high by a unit column driver circuit 220'.\nIn the next clock cycle T1, the PCSLD signal becomes high before the low-to-high transition of the PCSLE signal, so that the line CSL0 is deactivated. Next, after the second column address signals CA#1 (CA0-CA2) had reached the unit column pre-decoder circuit 180' and the time Tm1 has elapsed, when the CSL enable control clock signal PCSLE goes high again in response to the activation of the column address setting signal PYE. The unit column pre-decoder circuit 180' generates the second decoded column address signals DCA#1 (DCA0-DCA7). Here, assuming DCA1 signal is activated high, then a corresponding line CSL1 will be driven high by a unit column driver circuit 220'.\nThe other column select lines (such as CSL2 and CSL3) also will become activated and deactivated during next clock cycles (T2 and T3) in response to the other column address signals (such as CA#2 and CA#3) in the same manner as the above-mentioned.\nIn the above conventional memory device, the CSL enable control clock PCSLE should not go active until valid column address signals arrive at the pre-decoder 180 during each clock cycle Tc (where, c=0, 1, 2, . . . ). However, in the event the PCSLE signal goes high during a clock cycle (e.g., T1) before the valid column address signals CA#1 arrive at the column pre-decoder 180, owing to an insufficient delay time of Tm1, then the invalid column address signal CA#0 for the previous clock cycle T0 may be pre-decoded again by the corresponding column pre-decoder circuit 180' (refer to FIG. 3). Hence, the invalid decoded signal DCAB0 will be latched by the corresponding inverting latch 60 via the cascode inverter 40 (refer to FIG. 4). This leads to the activation of the column select line CSL0. Thereafter, when the valid decoded signal DCAB1 is activated by decoding the valid column address signals CA#1 and latched by a corresponding inverting latch 41 in cycle T1, the column select line CSL1 corresponding to the valid column address signals CA#1 also becomes active along with the invalid CSL0 line, causing an erroneous read/write operation. For the above reason, it is essential to ensure sufficient delay time Tm1 in the conventional memory device. This limits the memory access speed improvements.\nIn addition, according to the conventional memory device structure, a significant area penalty may result from the large reiterative layout area of the unit column driver circuit 220'.\nFurthermore, since the pull-up and switching transistors MP40a-MP47a and MP40b-MP47b within the respective cascode inverters 40-47 provide current leakage paths together with the inverters IV60b-IV67b of the inverting latches 60-67 during power-up, the conventional device has large power-up current dissipation."}
-{"text": "Semiconductor memory devices are typically classified into volatile memory devices and non-volatile memory devices. Volatile memory devices are subdivided into dynamic random-access memories (DRAMs) and static random access memories (SRAMs). Non-volatile memory types include erasable programmable read-only memories (EPROMS) and electrically erasable programmable read-only memories (EEPROMs). EEPROMs are increasingly used in system programming that requires continuous update or auxiliary memory devices. Particularly, flash EEPROMs are advantageous as mass storage devices because their integration density is high compared with conventional EEPROMs.\nFrequently, it would be convenient to be able to mix integrated circuit device types, such as EEPROMs and other memory devices, and bipolar integrated circuits, such as NPN transistors, onto a single integrated circuit chip. However, due to the inherently low breakdown voltage (approximately 10 volts) of typical wells used in BiCMOS technology and the need for a high programming voltage of an EEPROM memory device (approximately 14 volts), there has been no simple and economical way to integrate these two device types into a single integrated circuit. Previously, the problem has been avoided in the art by using additional masks to create high-voltage wells for EEPROMS."}
-{"text": "This invention relates to a new and distinct selection of Lagerstroemia indica, a member of the Lythraceae or Loosestrife family. Lagerstroemia indica cultivar Monink was discovered in a group of seedlings which originated from seed of Lagerstroemia indica `Little Chief`, seed was sown November 1983. During the summer of 1987, this new cultivar was selected at Monrovia Nurser Company, 18331 East Foothill Boulevard, Azusa, Calif. My new plant has been asexually reproduced by cuttings since 1987 at the above location in Azusa at Monrovia Nursery Company. The original seedlings displayed extreme variability, therefore the distinct phenotypic characteristics of my new selection that sets this plant apart from other Lagerstroemia indica plants would likely be lost through sexual reproduction. Therefore, sexual reproduction is prohibited and propagation is restricted to asexual reproduction by cuttings."}
-{"text": "1. Field\nThis relates to a dryer, and more particularly, to a dryer having enhanced dehumidifying power.\n2. Background\nIn a laundry treating apparatus having a drying function such as a washer or dryer, once washing and dehydration are completed, hot air may be supplied into the drum to evaporate moisture from the laundry, thereby drying the laundry. Such a dryer may include a drum rotatably provided within a cabinet, a drive motor to drive the drum, a blower fan to blow air into the drum, and a heating device to heat air conveyed into the drum. The heating device may use, for example, high-temperature electric resistance heat generated using electric resistance, or combustion heat generated by combusting gas."}
-{"text": "The present invention relates to a network relaying apparatus and a network relaying method, or in particular to a network relaying apparatus including a router of a computer network system which is capable of searching at high speed for a destination of a packet input and a network relaying search method.\nGenerally, in a network system, a network relaying apparatus such as a router or a bridge is used for connecting a plurality of networks. The router checks the destination address of a packet received from a network or a subnet connected, determines the destination of the packet, and transfers the packet to a network or a subnet which is connected with the destination router or host.\nFIG. 13 is a diagram showing a configuration of a conventional network relaying apparatus. In FIG. 13, a router 100 includes a routing manager (RM) 110, router buses 120, network interfaces (NIF) 130 and ports 140. Each port 140 is connected to an appropriate network 150.\nEach network interface 130 receives a packet from a network connected to the port 140, and transmits the received packet through the router bus 120 to the routing manager 110. The routing manager 110 includes a routing table for holding the routing information, and using this routing information, determines the network 150 of the destination from the address of the packet received, and transmits the packet to the network interface 130 of the port 140 connected to the network 150. The network interface 130 that has received the packet from the routing manager 110 sends out the packet to the destination network 150. The routing manager 110 updates and maintains the routing information held in the routing table based on the header information of the packet received, and has the function of overall management of the router 100.\nAn explanation will be given of the route search process for searching for a port outputting the next address to which the packet is to be transferred upon receipt of the packet and outputting the packet. Normally, the route search uses a route search table (routing table) prepared from the component definition information and the information obtained by exchange between the routers. The routing table is for searching the information (next hop information) as to the output port, the next hop address and whether the network is directly connected or not with a set of the network address and the network mask length as a key.\nAs another conventional system, JP-A-05-199230 (U.S. Pat. Ser. No. 5,434,863) discloses an internet-work system and a communication network system which can flexibly meet the size requirement of the network without adversely affecting the high-speed routing process. In these systems, a router manager and a plurality of routing accelerator modules are coupled to each other with a high-speed bus Also, each routing accelerator is connected with a plurality of independent communication ports. In these conventional systems, a plurality of the routing accelerators makes possible a high-speed routing and by adding the routing accelerators, the requirement for increasing the network size can be easily met.\nThe conventional router, however, cannot meet the requirement of the high speed lines such as the high-speed LAN (local area network) and the wide band ISDN (Integrated Services Digital Network) and ATM (asynchronous transfer mode) that have recently found applications. Also, the conventional router with only one routing means has the disadvantage that the number of ports and the communication traffic that can be supported are limited. It is therefore difficult to expand the configuration of the port menu of the router to a large size smoothly or to improve the performance in keeping with the port traffic volume.\nAlso, with the increase of internet users, the number of flows that the router is required to detect is on the increase. Therefore, it is necessary to set a multiplicity of flow conditions in the router. The increased traffic and the increased line speed of the internet, on the other hand, requires a shorter processing time per packet in the router. Even in the case where the number of set flow conditions is increased, therefore, the QoS (quality of service) control and the filter operation must be carried out at a high speed on the part of the router.\nIn setting flow conditions, on the other hand, a great variety of flow conditions set as desired by the router manager must be flexibly handled. The prior art fails to take this point into account.\nIn view of the above-mentioned points, an object of the present invention is to provide a network relaying apparatus and method for routing packets at high speed while assuring a high communication quality of service (QoS), a high reliability and security.\nAnother object of the invention is to provide a network relaying apparatus and method in which the flow conditions including the information for identifying users, the protocol information and the priority information can be set in great amounts, and in keeping with the increase in the line speed and the flow conditions, the flow can be detected at high speed so that the control operation for the communication quality including QoS control and filter can be realized.\nStill another object of the invention is to provide a network relaying apparatus and method in which the control operation flexibly meeting a great variety of flow conditions including the priority control, discard control and band control can be performed at high speed by improving the description of the flow conditions and the combination of the information including the source and the transfer destination.\nAccording to this invention, as described above, there are provided a network relaying apparatus and method for routing packets at high speed while at the same time assuring a high communication quality (QoS), a high reliability and tight security.\nAlso, according to this invention, the flow conditions including the information for identifying the users, the protocol information and the priority information can be set in great amounts in keeping with the increase in line speed and flow conditions, and the flow can be detected at high speed for realizing a high-speed QoS control and filtering. Further, the high-speed control operation flexibly meeting a great variety of flow conditions including the priority control, the discard control and the band control is made possible by improving the descriptiveness of the flow conditions and the combination of the information including the source and the transfer destination. Further, according to this invention, jobs of different categories (such as basic jobs and information jobs) can be combined into a single network.\nAccording to one aspect of the invention, there is provided a network relaying apparatus comprising:\nat least a network interface connected with at least a network;\nat least a routing processor including a packet buffer for storing input packets and a flow search table set separately for each of the input or output line number with action information corresponding to the information including the packet source and the packet transfer destination as an entry;\na routing manager for managing the internal components of the system; and\na connector for connecting the routing manager and each of a plurality of the routing processors;\nwherein the network interface outputs the input packet from the network to the routing processor; and\nwherein the routing processor includes means for storing the input packet from the network interface in a buffer memory, means for searching the transfer destination of the input packet stored in the packet buffer based on the stored header information, means for searching and reading only the entry corresponding to the input or output line number of the packet by referring to the flow search table, means for determining whether the information including the packet source and the packet transfer destination are coincident with the reference conditions in the entry read out, means for determining, in the case of coincidence, the control operation for the communication quality including the order of priority of packet transfer and the possibility of transfer in accordance with the action information in the entry, and means for outputting the input packet stored in the packet buffer and the output packet produced according to the header information, to the connector and the network interface.\nAccording to another aspect of the invention, there is provided a network relaying method for outputting the input packet input from a network to a transfer destination in a network relaying apparatus comprising at least a network interface connected to at least a network, at least a routing processor for routing the packet input from the network interface, a routing manager for managing the internal components of the system, and a connector for connecting the routing manager and each of a plurality of the routing processors;\nwherein the routing processor includes:\nmeans for setting the action information corresponding to the information including the packet source and the packet transfer destination as an entry separately for each input or output line number in a flow search table;\nmeans for storing the input packet in a buffer memory;\nmeans for searching for a transfer destination of the input packet stored in the packet buffer based on the stored header information;\nmeans for searching and reading only the entry corresponding to the input or output line number of the packet by referring to the flow search table;\nmeans for determining whether the information including the packet source and the packet transfer destination coincides with the reference conditions in the read entry;\nmeans for determining, in the case of a coincidence, the control operation for the communication quality including the priority of packet transfer or the possibility of transfer based on the action information in the entry; and\nmeans for outputting the input packet stored in the packet buffer and the output packet produced by the header information to the connector or the network interface."}
-{"text": "This invention relates to an improved oxidant and more in particular to an acidic, aqueous, oxidizing agent containing bromate and iodate ions.\nDying of various fabrics to impart a color to the fiber has been practiced for many centuries. The color must generally be permanently and uniformly distributed throughout the fiber and not merely superficially applied to the fiber as in painting. Many different types of natural and regenerated cellulosic fibers have been dyed to impart a color. For example, natural fibers, such as the vegetable fibers cotton, linen, jute, and flax have been dyed. Regenerated cellulosic fibers, such as viscose rayon and cellulose acetate, are those produced from natural materials which were altered by man to produce a desired textile material.\nIt has become accepted, and common, practice to color these materials with well-known sulfur and vat dyes. These dyes are water insoluble substances which are readily converted to a water soluble or leuco form by reducing the sulfur or vat dye in, for example, a solution containing an alkali and sodium sulfide or hydrosulfite.\nThe leuco forms of sulfur and vat dyes are water soluble and well known to be substantive to cellulosic fibers. After application to the fiber, the leuco dye must be oxidized to permanently color the fabric. The process of U.S. Pat. No. 3,775,047 oxidized the dye with an aqueous oxidizing solution including acetic acid and sodium or potassium iodate. U.S. Pat. No. 4,042,319 disclosed similar oxidation with an aqueous oxidant containing acetic or formic acid, an alkali bromate and an alkali iodate. Such oxidizing solutions are operable; however, it is desired to provide an improved material suitable to oxidize leuco forms of sulfur and vat dyes."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to a system and method for easy and rapid instruction of second language skills. More particularly, the present invention relates to a system and method for rapid instruction of second language skills that can accelerate profession-specific language learning skills.\n2. Related Art\nWhile it may be that the English language is the primary, or most used, language in the United States, the number of people in the U.S. who do not speak English is increasing daily. For instance, it has been estimated that the number of U.S. residents who speak Spanish may shortly outnumber the residents who speak English. This increase of non-English speaking residents posses difficulties for many businesses and organizations. For example, many businesses located in the U.S. are structured around the English language. Advertisements, menus, directions, etc. are often provided in English. Many of a business' employees may speak little or no foreign (that is, non-English) languages. Conversely, an increasing number of an employer's employees may have little or no skill in speaking English.\nIn order to provide effective service to people who do not speak English, a business currently has a limited number of options. First, it may recruit and hire employees who speak multiple languages and who are capable of communicating and serving speakers of languages other than English. This is problematic in that the number of available people in the job pool are correspondingly decreased and the cost of employing the employees is increased.\nSecond, a business may train its existing employees in other languages to enable them to better serve customers who speak languages other than English. This option is problematic in that current foreign language systems can be very costly and very time consuming, and aren't focused toward teaching the basic language skills needed to communicate on a basic level with a speaker of the second language. Also, most conventional training systems are focused primarily on teaching second language skills which encompass a large array of situations in which a person may require or use the second language, rather than focusing on specific needs for a particular profession.\nThese problems also arise in the instance where a business employs employees who speak little or no English, but who are nonetheless expected to serve English speaking customers."}
-{"text": "This invention relates to sensors that measure torque applied to a shaft. It particularly relates to sensors having a sleeve torsionally engaged with the shaft and sensing means responsive to torsional strain in the sleeve.\nMany known torque sensors operate by responding to magnetostrictive effects resulting from strain in a stressed member or transducer. Some of these are in commercial production. Efforts have been directed toward using magnetostrictive effects to measure the torque applied to the steering wheel by the driver of a motor vehicle. One known design torsionally engages a sleeve having desirable magnetostrictive properties to a portion of the steering wheel shaft. Another design uses the magnetostrictive properties of a current production steering wheel shaft to eliminate the cost of attaching a sleeve to the steering wheel shaft. In a third known design magnetostrictive material is beam or vapor deposited on a steering wheel shaft. The known designs have not proved entirely satisfactory. Known methods of attaching a sleeve require processes that are not easily adapted to large volume production. The same is true of beam or vapor deposition of magnetostrictive materials. Efforts to use the shaft itself have suffered from the difficulty of obtaining shafts consistently having desired magnetostrictive properties. For measuring steering torque in an automobile the ideal sensor would be inexpensive and compatible with existing steering wheel shafts.\nThe expression xe2x80x9ctorsionally engagedxe2x80x9d is used herein to describe engagement between a first element and a second element for transmitting torque therebetween. It includes engagement for transmitting torque by a rigid attachment such as a weld or adhesive joint or both elements being made of one piece of material. It also includes engagement by means that transmit only torque exemplified by a wrench socket engaging the head of a bolt. The expression xe2x80x9ctorsionally engagedxe2x80x9d is used to cover a broad range of torque transmitting engagement means that may or may not transmit forces in addition to torque.\nA torque sensor incorporating a sleeve of magnetostrictive material is described in U.S. Pat. No. 5,351,555 issued Oct. 4, 1994 to Garshelis. Particular attention is focused on the Garshelis patent because it is believed to offer the lowest cost sensor responsive to torque applied to a magnetostrictive sleeve. However, the invention is applicable to any torque sensor having a sleevelike transducer that is torsionally stressed when torque is applied to a shaft.\nThe Garshelis design provides a sleeve (xe2x80x9ctransducerxe2x80x9d) permanently magnetized in its circumferential direction. Garshelis discusses attachment of the transducer to the torsionally stressed shaft and (column 15 beginning at line 7) describes requirements which must be met by the chosen method of attachment:\nxe2x80x9cproper operation . . . requires that there be no slippage between any of the components at their interfaces. . . . Somewhat less obvious, but no less important, is the requirement that there be no inelastic strain in shaft 8 in any cross section which includes the transducer 4. Thus, all strains associated with the transmission of torque must be fully recoverable when the torque is relaxed.xe2x80x9d\nand in column 16 beginning at line 5\nxe2x80x9cAs already indicated, the transducer 4 and underlying shaft must act as a mechanical unit. Rigid attachment of the transducer 4 either directly or indirectly to shaft 8 is crucial to proper operationxe2x80x9d.\nIn fact, attachment by adhesive bonding (using known adhesives and known designs) or interference fit (Garshelis\"\" preferred method) do not satisfy the above quoted requirements. All known designs based on adhesive bonding or interference result in peak stresses exceeding the capabilities of the bond.\nIn column 16 beginning at line 5 and continuing through line 23 of column 17 Garshelis discusses three categories of torsional engagements between the transducer and the shaft. The categories are 1) salient point, i.e. splines, knurls, teeth etc. at the ends of the transducer mating with similar features on the shaft; 2) distributed, i.e. adhesive bonding or interference fit; 3) diffuse, i.e. welding or brazing the ends of the transducer to the shaft. The first xe2x80x9c1) salient pointxe2x80x9d and the last xe2x80x9c3) diffusexe2x80x9d work well but manufacturing methods for achieving these attachments are not easily adapted to automotive manufacturing procedures.\nAbout friction or adhesive bonding Garshelis states (column 16 lines 37 through 41):\nxe2x80x9cThis bonding limits the maximum measurable torque to a lower value than might otherwise be handled by the shaft 8 alone or transducer 4 alone, but is advantageous for other reasons as indicated previously.xe2x80x9d\nAccordingly, Garshelis expresses a known need for an xe2x80x9cadvantageousxe2x80x9d process such as adhesive bonding that does not limit the maximum measurable torque to xe2x80x9ca lower value than might otherwise be handled by the shaft 8 alone or transducer 4 alonexe2x80x9d. Garshelis goes on to state (column 16 lines 41 through 47):\nxe2x80x9cPress or shrink fits can be used to obtain the desired circular anisotropy, and can provide very substantial gripping forces which as a practical matter will not be broken by expected torques on shaft 8. With clean, degassed (and perhaps deoxidized) surfaces, the effective coefficient of friction can rise without limit and act somewhat like a weld.xe2x80x9d\nProviding xe2x80x9cclean, degassed (and perhaps deoxidized) surfacesxe2x80x9d on the elements before they are joined by press or shrink fits is expensive and time consuming. It is difficult to assure such qualities in many millions of parts as required for automotive production. It is not stated in the Garshelis patent but it is believed that to achieve in a press fit an effective coefficient of friction that xe2x80x9ccan rise without limit and act somewhat like a weldxe2x80x9d as stated in Garshelis the xe2x80x9cclean, degassed (and perhaps deoxidized) surfacesxe2x80x9d must be joined and heat treated at high temperatures in a suitable atmosphere for many hours. To obtain a shrink fit heat treatment is believed to be required both to achieve an effective coefficient of friction that xe2x80x9ccan rise without limit and act somewhat like a weldxe2x80x9d and to cause the shrinkage required for a shrink fit.\nAnother method of achieving an interference fit between the transducer and the shaft is described by Garshelis with reference to FIGS. 14, 15 and 16. In this method the shaft is hollow and an expander is drawn through the shaft to expand it thereby providing the desired hoop stress. This process also is believed to be difficult and expensive to implement in mass production of steering wheel shafts.\nThe following numerical examples will clarify the issues related to attaching a sleeve by adhesive or interference fit (without heat treatment or other processes to achieve an effective coefficient of friction that xe2x80x9ccan rise without limit and act somewhat like a weldxe2x80x9d). In column 10 lines 3 through 5 Garshelis cites the example of a shaft diameter of 0.5 inch (1.27 centimeters) and a transducer wall thickness in the 0.030 to 0.050 inch (0.076 centimeters to 0.127 centimeters) range. The wall thickness is important to achieve sufficient magnetic flux (Garshelis column 10 lines 24 through 31). From the well known fact that torque transmitted by a shaft is distributed as the third power of the radius it follows for the case of the aforementioned 0.5 inch diameter shaft that if the transducer and shaft have similar shear moduli (which is likely to be the case) 36 percent of the total torque will be transferred to the transducer in the case of 0.030 inch transducer wall thickness and 52 percent of the total torque will be transferred to the transducer in the case of 0.050 inch transducer wall thickness. A possible diameter of a steering wheel shaft of an automobile is 2 cm and it might be subjected to a maximum torque of 600 newton-meters (450 ft-lbs). Such a torque might be applied by a large healthy male driver after the wheel reached the end of its travel. At the one centimeter radius of the outer surface of the steering wheel shaft 600 newton-meters torque creates a tangential force of 60,000 newtons (13500 lbf). If 36 percent of the torque is transmitted to the transducer 21,600 newtons (4860 lbf) must be transferred between the transducer and the shaft by the attachment means. The fraction of the force transferred between the transducer and the shaft would be 36 percent in the case of the 2 centimeter diameter shaft if the inside diameter of the transducer is also 2 centimeters and the thickness of its wall is 1.2 millimeters (0.047 inches). The fraction would be much larger if the inside diameter of the transducer is larger and the thickness remains 1.2 millimeters. Assuming an adhesive shear strength of 10 newtons per square millimeter (1419 psi) and assuming means exist for providing constant shear stress over the area of adhesive attachment, transferring 21,600 newtons requires 21.6 square centimeters or 3.5 centimeters of shaft length of bonded area at each end of the transducer.\nThe second example is an interference fit. If the transducer wall thickness is 1.2 millimeters and is stressed to a hoop stress of 700 mpa (100,000 lbf/in2) and the coefficient of friction is 0.3, 1575 newtons (353 lbf) of shear force can be transferred per millimeter of length. Transferring the aforementioned 21,600 newtons of shear force requires about 1.4 centimeters of shaft length of contact with the shaft at each end of the transducer.\nIn summary, in the case of a two centimeter diameter steering wheel shaft, both bonding by adhesive and attachment by press fit would require contact with the shaft for one to four axial centimeters beyond each end of the active area of the transducer to transmit the forces encountered in operation assuming uniform shear forces. To prevent higher stresses that would cause adherence to fail the shear force must be distributed uniformly over the area of attachment. In fact, known technology does not enable the hereinabove reproduced requirements (Garshefis column 16 lines 37 through 41 and column 16 lines 41 through 47) to be achieved with any amount of adhesive or conventional press fit adherence area because the shear stresses peak at the ends of the attachment regions and exceed the maximum shear capabilities of adhesives and/or press fits.\nA substantial difference is now evident between xe2x80x9cdistributed attachmentxe2x80x9d (adhesive, friction) and xe2x80x9csalient point attachmentxe2x80x9d and xe2x80x9cdiffuse attachmentxe2x80x9d. In the latter two attachment is truly at the ends of the transducer and the transducer operates as a unit with the shaft. This is also true in the aforementioned case where the effective coefficient of friction rises without limit and acts somewhat like a weld which is believed to be properly categorized as a xe2x80x9cdiffuse attachmentxe2x80x9d. In the xe2x80x9cdistributed attachmentxe2x80x9d cases attachment forces are required to be distributed over lengths of shaft such as the aforementioned one to four centimeter attachment regions at each end of the transducer.\nIt will also be appreciated from the above numerical examples taken with the following that where the transducer has a constant thickness as illustrated in FIGS. 1, 3, 4 and 6 through 16 of Garshelis (all of the figures that illustrate transducers) the end portions of the transducer do not xe2x80x9cact as a mechanical unitxe2x80x9d with the steering wheel shaft unless the ends are effectively welded to the shaft. For the transducer to xe2x80x9cact as a mechanical unitxe2x80x9d with the steering wheel shaft it must twist as the steering wheel shaft twists over its entire length. However, if a constant thickness transducer is attached by adhesive or press fit, sufficient torque to twist the transducer and cause it to xe2x80x9cact as a mechanical unitxe2x80x9d with the steering wheel shaft is only achievable in a xe2x80x9ccentral regionxe2x80x9d between the aforementioned attachment regions. Outside the xe2x80x9ccentral regionxe2x80x9d the torque available to twist the transducer diminishes with distance from the xe2x80x9ccentral regionxe2x80x9d because the transmission of torque is xe2x80x9cdistributedxe2x80x9d and the twisting of the transducer diminishes as the torque diminishes with distance from the central regions causing the torsional strain of the transducer and the shaft to be different far from the central regions. In Garshelis\"\" words cited hereinabove: xe2x80x9cThis bonding limits the maximum measurable torque to a lower value than might otherwise be handled by the shaft 8 alone or transducer 4 alone.xe2x80x9d\nAn object of this invention is to provide a torque sensor transducer which can be attached by adhesive to a torque carrying shaft and which will then operate xe2x80x9cas onexe2x80x9d with the torque carrying shaft.\nA general object of this invention is to provide a torque sensor which also overcomes certain disadvantages of the prior art.\nThe present invention provides a torque sensor for measuring the torque applied to a shaft. It comprises a magnetostrictive sleeve torsionally engaged with two shear levelers. The shear levelers are bonded by adhesive to the shaft. The shear levelers have flared ends and regions of varying torsional elasticity that operate to level the shear stress in the adhesive. The term xe2x80x9clevelxe2x80x9d is used herein with reference to shear stress in adhesives to describe causing the shear stress to be constant and without peaks over the area bonded by adhesive. It may include being constant at two or more different levels at two or more areas bonded by adhesive.\nFurther, in accordance with the invention, the torque sensor is attached to the shaft by adhesive which is stressed in shear without stress peaks thereby enabling designs wherein the adhesive can transfer torques approaching the yield limit of the shaft.\nFurther, in accordance with the invention, the shear levelers have varying torsional stiffnesses to provide a uniform shear stress in the adhesive.\nFurther, in accordance with the invention, the shear levelers have flared ends and varying thickness adhesive at the flared ends further distributes stress in the adhesive and enables designs wherein the adhesive transmits torques that approach the yield torque of the shaft.\nFurther, in accordance with a first embodiment of the invention, low magnetic permeability isolation rings magnetically isolate the shear levelers from the magnetostrictive central segment.\nFurther, in accordance with the aforementioned first embodiment of the invention, the isolation rings are welded to the shear levelers and the magnetostrictive central segment.\nFurther, in accordance with the aforementioned first embodiment of the invention, the magnetostrictive central segment is pressed onto a stack of washers with crowned outer circumferences that maintains the transducer in its cylindrical shape and minimizes the torque that must be accumulated by the shear levelers. Great hoop stress in the magnetostrictive central segment is achieved by heat treatment after the magnetostrictive central segment is pressed onto the stack of washers.\nFurther, in accordance with the aforementioned first embodiment of the invention, each washer of the aforementioned stack of washers with crowned outer circumferences is coated with a thin layer of material that evaporates during heat treatment thereby leaving each washer separated from adjacent washers and therefore free to rotate without friction when the transducer is torsionally strained.\nFurther, in accordance with a second embodiment of the invention, the shear levelers are unitary with a low magnetic permeability middle segment upon which the magnetostrictive central segment is pressed and welded and possibly shrunk whereby great hoop stress in the magnetostrictive central segment is achieved which advantageously provides desirable magnetic properties.\nFurther, in accordance with the aforementioned second embodiment of the invention, the shear levelers are unitary with a low magnetic permeability middle segment upon which the magnetostrictive central segment is placed and welded and great hoop stress in the magnetostrictive central segment is achieved by expanding the middle segment and the magnetostrictive central segment together which advantageously provides desirable magnetic properties.\nFurther, in accordance with a third embodiment of the invention, the shear levelers are unitary with the magnetostrictive central segment and annular grooves are provided between the shear levelers and the magnetostrictive central segment. The grooves enhance magnetic anisotropy and provide surfaces against which force may be applied to facilitate installation of the transducer on the shaft.\nFurther, in accordance with the invention, a torque sensor comprises a circularly symmetric magnetic element centered on the rotation axis of a magnetostrictive element for providing a lower reluctance magnetic field path and less sensitivity to a bent shaft or other asymmetry.\nA complete understanding of this invention may be obtained from the description that follows taken with the accompanying drawings."}
-{"text": "Securing items during transport from one location to another location can be difficult for those without a vehicle. Depending on the distance from the retailer, the act of transporting the items can be physically demanding and frozen items can melt or defrost along the way."}
-{"text": "1. Statement of the Technical Field\nThe present invention relates to the field of markup language processing, and more particularly to the processing of frames in markup language.\n2. Description of the Related Art\nConventional markup can be visually presented through use of a content browser. Content browsers process display attributes embedded in markup to properly format content also contained within the markup. Notable variants of the content browser include the venerable Web browser, as well as the more recent extensible markup language (XML) browser. Regardless of the type of browser, all conventional markup processors are preconfigured to parse and interpret attribute tags embedded in markup. Examples of attribute tags include the well-known hypertext markup language (HTML) tags, <HEAD>, <BODY>, <H1>, <P>, <HREF>, <HTML> and <FRAMESET>.\nIn regard specifically to the <FRAMESET> attribute tag defining a set of displayable frames, many content distributors have incorporated frames in the design of Web pages, as frames allow Web sites to organize the presentation of disparate information in a logical and unified manner. Often, the enablement of multiple, scrollable regions of a single Web page forms the basis for the use of frames within the Web page. For example, it is known to use one frame to present a menu, while a second, adjacent frame can display the content associated with a particular menu choice. As another example, the results of a query can be presented in a first, scrollable frame, while the content associated with a particular result can be presented in an adjacent frame.\nWhile frames can provide the benefit of screen organization and scrolling, frames can pose some difficulty in the support of bidirectional environments. Bidirectional environments relate to the directional manner in which information can be presented in an electronic document. Specifically, western languages such as English present information visually from left to right. In contrast, middle-eastern languages such as Hebrew and Arabic present information visually from right to left. In the course of undertaking internationalizing the presentation interface of computer software, solutions have been proposed which address the problem of bidirectional text. For instance, in one solution, text can be rearranged in a display window according to the underlying language associated with the character codes of each word in the text.\nOther internationalization solutions address not only the underlying directional orientation of text, but also the arrangement of user interface controls in a graphical user interface. For example, in Kaplan, Internationalization with Visual Basic, (SAMS 199x), check boxes, labels, option buttons and text boxes can be horizontally \u201cflipped\u201d in terms of orientation to accommodate a bidirectional environment. The particular user interface control elements which can be flipped in the Kaplan reference, however, are flipped inasmuch as those control elements usually accompany text\u2014hence the need to re-orient the control. By comparison, the teachings of Kaplan explicitly inhibit the re-orientation of other user interface controls such as combo boxes, command buttons, scroll bars, list boxes, picture boxes and, most importantly, frames.\nWhile the use of frames in a conventional GUI as described in the Kaplan reference rightfully inhibits the re-orientation of frame elements because frames play little role in the GUI of a stand-alone application, the same cannot be said of an application whose interface relies upon the presentation attributes of a markup language. Specifically, frames play a crucial role in the presentation interface of a markup language defined user interface. Thus, in the context of markup, the directional implication of a language is not merely limited to the characters which form a word, or the words which form a sentence. Rather, the directional implication of a language can include the layout of the frames within the document itself. More particularly, in a right to left orientation, it can be preferable to horizontally re-orient adjacent frames in a content browser to accommodate a right-to-left environment.\nUnfortunately, existing content browsers and markup languages do not account for bidirectional documents. In fact, in the HTML specification, while the \u201cdir\u201d attribute can specify a directional orientation, including \u201crtl\u201d and \u201cltr\u201d\u2014right to left and left to right, respectively\u2014the HTML specification cannot account for a document whose orientation can vary from right to left and left to right. Rather, to accommodate the bidirectional circumstance, separate markup must be maintained for both cases of left to right and right to left configurations. As one skilled in the art will recognize, however, maintaining two sets of markup to support the presentation the same content in different directional orientations requires the maintenance and synchronization of both sets of markup\u2014a distinctly undesirable solution."}
-{"text": "1. Field of the invention\nThe present invention relates to landscape edging used to demarcate areas of lawns, gardens and the like, and more particularly to an illuminated landscape edging.\n2. Description of the Prior Art\nLandscape edging is very widely used to provide demarcation between areas of lawns, gardens and the like. For instance, landscape edging may be used to form a barrier between botanical types, wherein, for example, on one side thereof is a lawn and on the other side thereof is a flower garden. In order to provide a barrier to passage therethrough of flora, landscape edging is implanted into the ground along the line of selected demarcation, and a portion thereof is located above the surface of the ground.\nOne very common type of landscape edging is depicted in FIG. 1. This prior art landscape edging 10 is characterized by a planar member 12 which is structured to be implanted into the ground and by a tubular member 14 which is structured to be located above ground. The planar member 12 is provided with a medial ribbing 16 and usually provided with an end hook 18. Both the medial ribbing 16 and the end hook 18 serve to anchor the planar member into the ground. This prior art landscape edging 10 is constructed of a plastic material, which is believed to be invariably black and opaque. The prior art landscape edging is elongated and is supplied in long rolls which may be cut into sections to suit the length needs of a particular landscaping job. Sections of the prior art landscape edging are serially joined by a tubular connector 20 inserting into open ends of adjoining tubular members 14.\nIncreasingly, homeowners are installing landscape lighting to illuminate their lawns and other landscape features. Usually, this form of illumination is in the form of a series of discrete lighting units which are partly implanted into the ground and which operate on a low voltage line supplied from the home utility service. These units can sometimes be the target of vandals or may be the accidental victims of an untoward incident, such as being hit by a lawn mower.\nWhat is needed in the art is to somehow combine the benefits of landscape edging with the attractiveness of landscape lighting."}
-{"text": "1. Field of the Invention\nThe invention relates to degaussing coils or degaussers designed to be fitted into cathode-ray tubes for color televisions. The degaussing coil is constituted by a winding of enamelled electrical wire shielded by a heat-sealed insulator sheath. This wire is fastened to the glass fixtures of the mask picture tube and it is connected to the electrical circuit of the television set by a pair of conductors that end in a bipolar connector. The conductors are electrically connected to the two ends of the winding in a region that is not covered by the insulator sheath. This uncovered region is then confined in an insulating envelope that ensures the electrical shielding and the mechanical strength of the connected parts. The envelope can be formed by taping or by two shells of synthetic organic resin forming an interconnection case.\n2. Discussion of the Background\nCoils manufactured according to taping techniques or by the positioning of a case require equipment whose difficulty of handling entail considerable extra costs. These coils may have a defect in the sealing quality at the interconnection and they lack aesthetic quality."}
-{"text": "A lens array optical system including a plurality of lens portions is used for an imaging lens for a compound-eye imaging apparatus, a secondary image forming lens in an autofocus module of a single-lens reflex camera, an illumination lens in a semiconductor exposure apparatus, a condenser lens in a liquid crystal projector panel, and the like. The lens array optical system can be manufactured at low cost with optical plastic having high processability. However, the optical plastic not only has a high thermal expansion coefficient bus also is inferior in durability and light transparency, and accordingly is not suitable to be used stably on a high-temperature or high-humidity severe condition. Therefore, optical glass that has a low thermal expansion coefficient and is superior in durability under high temperature or high humidity is required to be used for the lens array optical system in order to enable stable use in various environments.\nAs a method for manufacturing a lens array optical system made of optical glass having multiple spherical or aspheric lenses, a method is proposed in which preforms made of glass with a larger curvature than a concave curvature of a forming die are placed, one by one, in transfer portions and are hot stamped (refer to Patent Literature 1). An excess of the preforms over the volumes of cavities for forming optical surfaces of the multiple lenses flows into junctions of the cavities and are integrally fused. However, there is a problem with the method of Patent Literature 1 that it takes much time and trouble to place the preforms on the optical transfer surfaces on the die. Moreover, air may remain in the junction (fused portion or joint) due to variations in the preforms to result in a reduction in the strength of the junction of the preforms. Moreover, there is also a problem that as the number of lens portions increases, the forming process becomes more difficult, and accordingly it becomes impossible to obtain a lens array having an intended optical surface shape.\nMoreover, as another method for manufacturing a lens array optical system, a method is also proposed which does not use a die, supplies a liquid material to a substrate where a plurality of through-holes is formed, and forms a plurality of lenses on the substrate (refer to Patent Literature 2). However, there is a problem with the method of Patent Literature 2 that it is difficult to form an optical surface of the lens into an aspheric shape so that it is not possible to appropriately design the optical surface of the lens portion and optically correct various aberrations.\nMoreover, as another method for manufacturing an optical glass lens array optical system, the drop method is known which forms a lens by dropping molten glass into a die. The drop method has few constraints on the shape of an optical surface since the forming process is performed in a state where the viscosity of the glass is relatively low. Moreover, a reduction in the occurrence of forming failure ascribable to the fluidity of a material can be expected. As the method for manufacturing a lens array by the drop method, Patent Literature 3 describes a method for manufacturing a lens array having 2\u00d72 lens portions by dropping a glass drop on the center of one of dies having 2\u00d72 optical transfer surfaces and pressing it with the other die. However, Patent Literature 3 does not describe the production of a lens array having lens portions more than 2\u00d72, such as 4\u00d74 or 5\u00d75, in a lattice form. If an attempt is made to obtain a lens array having 4\u00d74 lens portions by dropping a glass drop on a center portion of, for example, 4\u00d74 optical transfer surfaces in accordance with the manufacturing method described in Patent Literature 3, an optical transfer surface that has been insufficiently filled with the glass may be produced as described below. Moreover, in Patent Literature 3, a projection is provided to change the flow of the glass drop to make adjustments such that the glass drop flows along the transfer surface of an optical surface near its edge side close to the projection, the edge side having a large inclination angle. Accordingly, an attempt is made to transfer the optical surface with high accuracy. However, there are problems that the manufacturing cost may increase if a glass projection is provided to the die as a preliminary process, and the die processing cost is high and the life of the die is influenced if the projection is formed upon die processing. If an attempt is made to drop a plurality of glass drops individually and respectively onto the optical transfer surfaces, forming failure may be invited, or gas entrainment may occur upon reaching the die, due to time differences in die arrival timings among the glass drops. Moreover, a problem also arises in which it becomes difficult to perform a forming process by pressing to a desired thickness since it becomes easy to become cold and harden as the result of a reduction in the volumes of the individual glass drops."}
-{"text": "A known voltage converter and a known inverter have been used when driving a three-phase motor. The voltage converter boosts an input voltage to a predetermined voltage. The inverter changes a frequency of an output of the voltage converter. Meanwhile, plural switching elements are provided at the voltage converter and the inverter. The switching elements self-heat by energization, and electrical characteristics and a rating are specified in response to an ambient temperature level of the switching elements. Here, the ambient temperature is desired to be considered upon the use of the switching elements. The technology in which the switching elements are used under the consideration of the ambient temperature is disclosed in JP2013-48515A (hereinafter referred to as Patent reference 1).\nAn electric automobile disclosed in Patent reference 1 is provided with an inverter supplying electric power to a three-phase motor. The inverter includes plural switching elements. The electric automobile is provided with a temperature sensor, a current sensor, and a voltage sensor. The temperature sensor measures a temperature level of a refrigerant cooling the switching elements. The current sensor measures an output current of the inverter. The voltage sensor measures an input voltage inputted to the inverter. A temperature correction value is calculated based on a measurement data of the current sensor per switching element, a measurement data of the voltage sensor per switching element, and a duty ratio of the switching element. The temperature correction value is added to a measurement result of the temperature level of the refrigerant to estimate the temperature level of the switching elements.\nAccording to a device disclosed in Patent reference 1, detection results of the plural sensors are used to estimate the temperature level of the switching elements. In this case, because the detection results include errors of measurement, the errors of measurement included in the estimate results of the temperature levels of the switching elements may be increased. Accordingly, when the switching elements are energized, the electric characteristics and the rating are desired to include margins. Accordingly, the device disclosed in Patent reference 1 may not be able to fully use a capability of the switching elements.\nA need thus exists for an energization control system and a sensor unit which is not susceptible to the drawback mentioned above."}
-{"text": "1. Technical Field\nThe invention relates to speed transmission systems for automobiles, trucks, and other machines requiring the input shaft speeds to be different from the output shaft speeds, and in particular positive variable speed transmissions.\n2. Background Art\nThe art useful for the understanding, searching and examination of the invention is that of mechanisms and machine elements in general and more particularly positive variable speed changing mechanisms as classified in Class 74 subclasses 217,230,230.21, 244, 681, and 689.\n3. Disclosure of the Invention\nThe invention improves the conventional speed transmission systems by the introduction of unique moving bevelled mating plates between input and output shafts. The plates, named as \"Ruff Plates\", have unique surfaces that determines the rotational speed of one shaft to the other by the positioning of a plurality of continuous series of spaced-apart interconnected balls between the plates. The plates are thicker in the center and decrease to practically zero thickness at the peripheral edges. The mating surfaces of the plates comprises concentric series of crests and valleys increasing from the center to the peripheral edges in aligned sectors that allow the balls to pass between."}
-{"text": "As is well known, various tool holders have been utilized in the prior art which interface with a rotating spindle of a machine such as a milling or boring machine to securely interface a cutting tool to the machine during the cutting of a work piece. A rotary cutting machine typically includes a motor-rotated spindle to which a tool holder is attached, the tool holder being configured to accommodate a shank portion of a cutting tool which is ultimately used to cut a work piece. The attachment of the tool holder to the spindle is generally accomplished by providing a cavity in the spindle into which an upper end or shank portion of the tool holder is secured, as with an externally threaded bolt which is advanced through a portion of the spindle and is threadably received into an internally threaded bore extending axially within the shank portion of the tool holder. In most prior art tool holders, a central aperture is also formed in a lower end or mounting portion of the tool holder for receiving the shank portion of the cutting tool which is to be interfaced to the milling or other machine via the tool holder. Subsequent to the insertion of the shank portion of the cutting tool into the central aperture of the tool holder, the shank portion of the tool holder is drawn or pulled tightly into the spindle by the tightening of the above-described bolt so as to rigidly maintain the cutting tool within the tool holder.\nThough interfacing the cutting tool to the cutting machine, prior art tool holders typically possess certain deficiencies which detract from their overall utility. More particularly, while slower rotational speeds generally permit the cutting machine to perform adequately, high speed cutting, which is preferred, can cause substantial difficulty in producing a satisfactory work as a result of the development of vibratory forces that occur between the spindle and the tool holder. Specifically, at higher rotational cutting speeds, the cutting tool often begins to chatter or vibrate such that adequate tool control cannot be maintained and damage to the work piece, along with potential injury to the operator, can realistically occur. More particularly, the transfer of the harmonic resonance into the tool holder may give rise to slight movements thereof relative to the spindle, and in extreme cases may result in the tool holder loosening within the spindle. As will be recognized, the resonance of the tool holder relative to the spindle results in the cut in the work piece being substantially out of tolerance.\nA further deficiency with prior art tool holders is that the manner in which the shank portion of the cutting tool is secured within the central aperture of the tool holder often results in the non-concentric mounting of the cutting tool within the tool holder. Such non-concentric mounting is unacceptable in modem, high tolerance machining applications such as those performed on a vertical milling machine wherein minor variations in the concentricity of the cutting tool within the tool holder often cause catastrophic failure in the cutting operation.\nThe present invention addresses these and other deficiencies of the prior art tool holders by providing a tool holder which includes a dampening member for eliminating the harmonic resonance which typically occurs during the use of cutting tools in high speed milling applications. In the tool holder constructed in accordance with the present invention, the non-concentric mounting of the shank portion of the cutting tool within the tool holder is also substantially eliminated by the heat shrinking of the shank portion of the cutting tool within the tool holder. These, as well as other features and advantages attendant to the present invention will be discussed in more detail below."}
-{"text": "1. Field of the Invention\nThe invention pertains to fire resistant firesleeve assemblies as used with hose or pipe systems, and particularly pertains to the sealing of the cut end of a firesleeve utilizing fibrous heat insulative components.\n2. Description of the Related Art\nHose systems, particularly fuel, lubricating and hydraulic hose used within aircraft engine compartments, often require fire resistant protection to minimize the likelihood of fire or high temperatures damaging the flexible elastomeric hose of the hydraulic circuits. Typically, resistance to fire and high temperature is achieved by encompassing the hose and portions of the associated hose fittings with a protective firesleeve. The firesleeve often consists of a silicone rubber tube surrounding the hose and/or fitting and the firesleeve usually includes an inner insulative material such as asbestos or glass fiber to provide insulation against the transfer of heat.\nTypically, firesleeve material is shipped and stored in indeterminate lengths, often wound upon reels, and the firesleeve is cut to the desired length to accommodate the length of hose with which it is to be used. As the cutting of the firesleeve to the desired length will expose the end of the firesleeve and the fibrous material thereof it is required that the end of the firesleeve be treated or sealed so as to discourage contamination of the firesleeve inner fibrous material and prevent the \"wicking\" of moisture, fluids, or other liquids into the fibrous material by capillary action.\nTo \"seal\" the cut ends of firesleeves, the conventional practice is to dip the cut firesleeve end into a silicone formulated \"end-dip\" liquid sealant so as to saturate the braided fibrous material adjacent the firesleeve end. Thereupon, the silicon end dip composition is dried and processed prior to the firesleeve being installed upon the associated hose. Such processing of the ends of firesleeves often requires a 24 hour delay, and the handling of a dipped firesleeve is messy, produces dripping, and such end dipping is hazardous and may not produce a liquid impervious end seal.\nFurther, the silicone formulated end dip used to seal firesleeves may be classified as a hazardous material requiring special and expensive transportation procedures. Further, end dipped firesleeves do not produce a significant frictional engagement with the hose or hose fitting unless the clamping pressures are unusually high and there is a tendency for the firesleeve to axially slide on the hose, possibly exposing a portion of the hose adjacent the fitting to fire in high temperatures.\n3. Objects of the Invention\nIt is an object of the invention to provide a fluid impervious end seal for hose or pipe firesleeve assemblies which is inexpensive, requires no dipping or application of fluids, requires no hazardous materials, and eliminates assembly and processing delays.\nAnother object of the invention is to provide a seal for the cut ends of hose or pipe firesleeves wherein the seal is formed by a flexible resilient liquid impervious cuff enclosing the firesleeve cut end which totally confines the cut end to produce an effective liquid seal, enhances fire protection, and produces an attractive aesthetically finished hose assembly."}
-{"text": "An apparatus 22 for double-side polishing semiconductor wafers as shown in FIGS. 4 to 7, wherein only the main functional members are shown in FIG. 4, is well known as an apparatus for polishing the two sides of each of semiconductor wafers W at the same time.\nThe apparatus 22 shown in FIG. 4 comprises a lower polishing turn table 24 and upper polishing turn table 26 both of which are arranged in a face to face relation with each other. Lower and upper polishing pads 24a, 26a are fixedly applied on the upper face of the lower polishing turn table 24 and the lower face of the upper polishing turn table 26 respectively. The lower and upper polishing turn tables 24, 26 are respectively rotated in opposed directions by a driving means which is not shown for clarity.\nA sun gear 28 is at the center of the lower polishing turn table 24 and held on the surface thereof and an internal gear 30 in an annular shape is located near the outside of the periphery thereof. Each of wafer carriers 32 has the shape of a disc and lies between the upper surface of the lower polishing pad 24a and the lower surface of the upper polishing pad 26a. The wafer carriers 32 are moved around in a space between the lower and upper polishing pads 24a, 26a, while each of the wafer carriers 32 turns not only about its own center but revolves around the centers of the polishing pads 24a, 26a by the actions of the sun gear 28 and internal gear 30. Each of the carriers 32 has a plurality of carrier holes 34 and semiconductor wafers W to be polished are placed in the respective carrier holes 34. In polishing the wafers W, polishing slurry is supplied between a wafer W and each of the polishing pads 24a, 26a through a hole 38 communicating between the upper and lower surfaces of the upper polishing turn table 26 from a nozzle 36. As a wafer carrier 32 turns about its center and at the same time revolves around the centers of the lower and upper polishing turn tables 24, 26, the wafers W placed in the carrier holes 34 turn about the respective centers of their own and at the same time revolve around not only the centers of the polishing turn tables 24, 26 but also the center of the carrier 32, so that the wafers W are moved around between the lower and upper polishing pads 24a, 26a and the two sides of each wafer W are polished at the same time.\nThe inner peripheral edge of a carrier hole 34 has an annular member A secured along the same edge for the purpose of preventing cracking or chipping of the edge of a semiconductor wafer W. The annular member A is produced by injection molding of plastics or it may be produced from glass fiber reinforced epoxy resin composite. Structures similar to or same as the wafer carrier 32 including the annular member A are respectively disclosed in laid open publications of Japanese Utility Model Application No. 57-177641, 58-4349 and 58-59559, and laid open publication of Japanese Patent Application No. 6-226618.\nWhen double-side polishing of a semiconductor wafer W is conducted by the use of the conventional apparatus 22 for double-side polishing semiconductor wafers above mentioned, a surface of a peripheral portion of an edge of the wafer W can be polished by the function of an annular member A as a polishing buff because of a relative rotational speed of the wafer W to the carrier hole 34 due to the rotation of the wafer W within the carrier hole 34, but the upper and lower chamfered portion of the edge of the wafer W remain unpolished. Therefore, generation of particles from such portion of the edge of the wafer W cannot be prevented in the conventional double-side polishing operation. In light of the problem, it is required that the whole edge portions are polished in double-side polishing. There are additional other problems in realizing the perfect polishing of the whole edge portions in double-side polishing that the perfect edge polishing as an additional process step is a costly operation because of its technical difficulty and moreover, the main surfaces of the wafer W are contaminated during edge polishing process."}
-{"text": "Zinc, aluminium and/or combinations of aluminium and zinc (Al/Zn), are widely used as surface coatings, particularly but not exclusively for steel for protection against rust and corrosion. In practice, however, the zinc or Al/Zn coatings are susceptible to white rust or black rust respectively when exposed to the atmosphere due to reactions with moisture. Such rust is detrimental to the surface and generally makes coated steel substrates unsaleable despite the fact that the overall service life of the coated steel may remain the same and the formation of rust generally interferes with finishing operations. The ability to resist such corrosion is referred to herein as wet stack performance.\nIn order to inhibit the formation of rust on coated surfaces it is generally accepted that the treatment of a surface with a chromate imparts anti corrosive properties and this type of treatment is generally referred to as chromate passivation. However chromate is highly toxic to exposed workers and due to its high toxicity, disposal of chromium residues is difficult. Further in various markets the yellow discolouration of the treated coated surfaces is considered to be an unacceptable product attribute.\nIn order to overcome the problems associated with chromate passivation, phosphate coatings have been used. However the anticorrosion properties of phosphate have been found to be far inferior to the above mentioned chromate treatment.\nU.K. Patent Application No. 2,070,073 discloses an anti corrosive treatment for preventing white rust on galvanized steel comprising applying to the surface of a galvanized steel sheet an acidic solution containing molybdic acid or a molybdate in a concentration of 10-200 g/l calculated as molybdenum and adjusted to a pH of 1 to 6 by addition of an organic or inorganic acid. However, it has been found that the corrosion resistance of aluminium and aluminium alloy surfaces treated with the above solution is inferior to the chromate treated substrates under certain conditions and the treated surfaces suffer from an undesirable degree of discolouration. Moreover modybdate treated surfaces have been shown to change from a pale yellow/blue to a strong green colour when stored for periods of time in excess of 24 hours.\nAccordingly it is an object of the present invention to provide means for avoiding and/or ameliorating at least some of the above discussed disadvantages of the prior art."}
-{"text": "1. Field of the Invention:\nThe invention relates to an integrated circuit for detecting a received signal, which can be used, for example, in radio receivers and in particular in mobile radio receivers, as well as a circuit configuration with an integrated circuit.\nVoltage-controlled oscillators (VCO) which are tuned with a phase locked loop (PLL) have become established in radio reception technology. The oscillator frequency of the voltage-controlled oscillator is tuned using the phase locked loop with a specific step size, for example with a frequency step size of 100 KHz in frequency-modulated radio in Europe. During the transmitter search run, the entire FM band (87.5 MHz to 108 MHz) is preferably scanned in 100 KHz increments. The oscillator frequency of the voltage-controlled oscillator fOSC is typically 10.7 MHz above the input frequency to be received, that is to say between 98.2 MHz and 118.7 MHz.\nDuring the transmitter search run, as far as possible all the transmitters which are worth receiving are to be found and their input frequency or the corresponding oscillator frequency fOSC of the voltage-controlled oscillator are to be stored. If a transmitter is detected at any particular reception frequency, the phase locked loop is stopped at the a instantaneous value by the microcontroller. The corresponding values are stored in the microcontroller so that these transmitter stations which are determined can be set directly at a later time. In order to be able to detect the transmission stations, one orxe2x80x94for sake of increased detection precisionxe2x80x94a plurality of criteria are required in order to stop the voltage-controlled oscillator at the corresponding oscillator frequency. It is therefore necessary to use one, or preferably a plurality of, evaluation criteria in the receiver in order to to detect unambiguously the presence of an input signal which is worth receiving.\nIn order to increase the reliability with which a transmitter can be detected unambiguously, it is possible to use not only the field-strength signal (level of the received signal) but also the multipath signal which indicates whether the received signal is scattered by multipath reception. To do this, the field-strength signal and the multipath signal are fed to a microprocessor which digitizes these signals and evaluates them through the use of criteria which are predefined in the microprocessor. Finally, in the microprocessor it is decided whether the transmitter is one which is worth receiving. The disadvantage of such a circuit configuration is that a plurality of lines are required from the receiver module to the microprocessor and unavoidable data traffic occurs on the data bus during the transmitter search run. However, this not only unnecessarily severely loads the microprocessor but the data traffic also constitutes a permanent source of interference in the sensitive reception system. In order to keep the data traffic as low as possible, the transmitter search criterion of intermediate frequency counting is therefore frequently not included, which however has the disadvantage of less precise transmitter identification. In many cases, for a more reliable search run stop it is not sufficient to evaluate only one of the search run stop criteria including the field-strength signal, multipath signal, zero crossover of the S curve and the limited intermediate frequency transmitted from the input mixer to the intermediate frequency amplifier. It is difficult to reliably detect weakly receivable transmitters only through the use of the field-strength.\nEuropean Patent Application No. EP 0 430 469 A2 discloses a circuit for detecting a received signal for FM receivers. The detector contains an intermediate frequency detector, a field-strength comparator and a noise signal detector whose output signals are combined with one another in a logic switching element to form a signal path. A microprocessor is driven as a function of this.\nPatent Abstracts of Japan, Volume 009, No. 248 (E-347), Oct. 4, 1985 (JP-A-60-096913) discloses how to use a noise detector to detect noise on the basis of multipath signals.\nEuropean Patent Application No. EP 0 335 141 A2 discloses a configuration for detecting various reception characteristic variables including a multipath reception, the outputs of which are coupled to a logic element.\nIt is accordingly an object of the invention to provide an integrated circuit for detecting a received signal which overcomes the above-mentioned disadvantages of the heretofore-known circuits of this general type and which permits a transmitter identification as quickly and reliably as possible while avoiding interference.\nWith the foregoing and other objects in view there is provided, in accordance with the invention, in combination with a microprocessor, an integrated circuit for detecting a received signal, including:\nan intermediate frequency detector configured to detect an intermediate frequency and supplying a first search run stop signal if the intermediate frequency is within a given range;\na field-strength comparator configured to supply a second search run stop signal if a field-strength of the received signal exceeds a field-strength setpoint value;\na multipath comparator configured to supply a third search run stop signal if a multipath signal exceeds a given multipath setpoint valve;\na logic component operatively connected to the intermediate frequency detector, the field-strength comparator, and the multipath comparator, the logic component logically combining the first, second and third search run stop signals with one another and forming a binary stop signal, the logic component having an output and providing the binary stop signal as a statically present signal at the output of the logic component, and the microprocessor receiving, as an input signal, the binary stop signal provided at the output of the logic component;\na first analog/digital converter connected upstream of the field-strength comparator, the first analog/digital converter digitizing a field-strength signal;\na first serial/parallel converter connected between the first analog/digital converter and the field-strength comparator;\na second analog/digital converter connected upstream of the multipath comparator, the second analog/digital converter digitizing the multipath signal; and\na second serial/parallel converter connected between the second analog/digital converter and the multipath comparator.\nIn other words, the integrated circuit according to the invention for detecting a received signal has an intermediate frequency detector which supplies a first search run stop signal if the intermediate frequency lies within a specific range. Furthermore, the integrated circuit has a field-strength comparator which supplies a second search run stop signal if the field-strength of the received signal exceeds a field-strength setpoint value. The invention additionally has a multipath comparator which supplies a third search run stop signal if the multipath signal exceeds a specific multipath setpoint value. Furthermore, a logic component is provided which logically combines the three search run stop signals with one another to form a binary stop signal which is statically present at the output of the logic component and is made available to a microprocessor as its input signal.\nThe data traffic between the microprocessor and the integrated circuit is reduced, which results in a severe reduction in interference.\nThe integrated circuit has the advantage that the overall space required for the radio receiver can be reduced as a result of the integration of a first analog/digital converter, which is connected upstream of the field-strength comparator and serves to digitize the field-strength signal, into the integrated circuit. Likewise, the external electrical lines, which therefore run outside the integrated circuit, can therefore be reduced, which also reduces the susceptibility to interference.\nA second analog/digital converter which is connected upstream of the multipath comparator and which serves to digitize the multipath signal also entails the aforementioned advantages.\nThe integrated circuit has the advantage that accelerated signal processing is possible in the field-strength comparator by virtue of a first serial/parallel converter which is connected between the first analog/digital converter and the field-strength comparator.\nThis advantage also applies if a second serial/parallel converter is connected between the second analog/digital converter and the multipath comparator.\nAccording to another feature of the invention, the integrated circuit can advantageously be adapted to the ambient conditions by setting the field-strength setpoint value and the multipath setpoint value to lower values in areas or conditions with weak reception than in areas or conditions with strong reception, so that rapid transmitter signal detection is still possible.\nIn an advantageous embodiment of the integrated circuit according to the invention, the field-strength setpoint value, the multipath setpoint value and the range within which the intermediate frequency is to lie are determined by the microprocessor. This permits a flexible, rapid and simple adaptation of the setpoint values to the ambient conditions without the computing capacity of the microprocessor being appreciably restricted thereby.\nAccording to another feature of the invention, a third serial/parallel converter having an output side coupled to the field-strength comparator and to the multipath comparator is provided, the third serial/parallel converter receiving the field-strength setpoint value and the multipath setpoint value in a serial manner.\nWith the objects of the invention in view there is also provided, a circuit configuration, including:\na microprocessor having a microprocessor input terminal and at least one microprocessor output terminal, the microprocessor input terminal having a single line; and\na detection circuit including:\nat least one detection circuit input terminal;\na detection circuit output terminal connected to the single line of the microprocessor input terminal;\nan intermediate frequency detector configured to detect an intermediate frequency and supplying a first search run stop signal if the intermediate frequency is within a given range;\na field-strength comparator configured to supply a second search run stop signal if a field-strength of a received signal exceeds a field-strength setpoint value;\na multipath comparator configured to supply a third search run stop signal if a multipath signal exceeds a given multipath setpoint valve;\na logic component operatively connected to the intermediate frequency detector, the field-strength comparator, and the multipath comparator, the logic component logically combining the first, second and third search run stop signals with one another and forming a binary stop signal as a statically present output signal, the logic component providing the binary stop signal to the detection circuit output terminal, the at least one microprocessor output terminal being connected to the at least one detection circuit input terminal for transmitting the field-strength setpoint value and the multipath setpoint value as a serial data word, and the given range within which the intermediate frequency is to be;\na first analog/digital converter connected upstream of the field-strength comparator, the first analog/digital converter digitizing a field-strength signal;\na first serial/parallel converter connected between the first analog/digital converter and the field-strength comparator;\na second analog/digital converter connected upstream of the multipath comparator, the second analog/digital converter digitizing the multipath signal; and\na second serial/parallel converter connected between the second analog/digital converter and the multipath comparator.\nIn other words, a circuit configuration according to the invention includes an integrated circuit and a microprocessor in which the integrated circuit has an output terminal for the binary stop signal which is connected to an input terminal of the microprocessor which includes a single line and in which the microprocessor has at least one output terminal which is connected to at least one input terminal of the integrated circuit in order to transmit over it the field-strength setpoint value and the multipath signal value as a serial data word, and the range within which the intermediate frequency is to lie.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in an integrated circuit for detecting a received signal and a corresponding circuit configuration, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawing."}
-{"text": "1. Field of the Invention\nThe present invention relates to a system for collecting scraps produced when various materials such as metals and plastics are processed by processing machines such as pressing machines and lathes and, more particularly, relates to a system which can efficiently collect scraps by the use of a suction force generated by a flow of high-pressure air.\n2. Statement of the Prior Art\nHeretofore, two typical systems have been used to collect scraps produced from a pressing machine, i.e. a press, which is one example of a machine for processing materials, when punching is carried out thereby.\nIn the first system, scraps are collected by simply receiving them as they fall due to their own weight. In this case, a box for collecting scraps is placed below the punching tool of the pressing machine, and scraps produced by punching of sheets fall due to their own weight and are received in the box and collected.\nThe second system is a dust collector. One example of such a dust collector system is disclosed in, e.g., Japanese Utility Model Laid-Open No. 61-24157. This system is designed as a dust collector for a portable air grinder, and is characterized by having an air ejector fixedly provided on a lid member of a cylindrical tank, an air filter provided in the tank at the downstream end of a suction pipe leading to the air ejector, a compressed-air inlet passage in the ejector connected with an air supply pipe through a valve, said lid member being provided with a dust collection opening, and a dust suction duct for connecting an abrasive wheel cover of the grinder to the dust collection opening. The air supply line to the grinder is connected by a T-joint to the compressed air inlet passage.\nHowever, such conventional scrap collection systems have the following disadvantages. A problem with the first system which simply collects falling scraps is that scraps are often prevented from falling because instead they are unintentionally deposited in a mold or the like, and are responsible for making dents in the products molded thereby, thus resulting in an increase in the number of defective products, a shortening in the service life of tools and a lowering of the efficiency of the system.\nA problem with the second system, namely the dust collector, is that the size of the duct collector becomes very large, so that not only is there a need to provide a space sufficient for its installation, but the complicated equipment also leads to an increased cost. Another problem is that as the sharpness of the processing tools deteriorates, scraps are likely to ascend while sticking to the punches because of a lack of suction force of the dust collector, or holes or openings are clogged by scraps."}
-{"text": "1. Field of the Invention\nThis invention relates to a novel compound and a chiral smectic C (SC*) liquid crystal composition containing the compound and being useful for liquid crystal display elements.\n2. Description of the Related Art\nAt present, as to liquid crystal display elements, TN (Twisted Nematic) type display mode has been most broadly employed, but as far as the response speed is concerned, such TN type display elements are inferior to emissive type display elements (such as those of electroluminescence, plasma display, etc.). Although various improvements in this respect have been attempted, it appears that improvement to a large extent has not yet been realized. Thus, various liquid crystal display devices based on a different principle from that of TN type display elements have been attempted. As one of such devices, there is a display mode utilizing a ferroelectric liquid crystal (N. A. Clark et al: Applied Phys. lett., 36, 899 (1980)). This mode utilizes the chiral smectic C phase (hereinafter abbreviated to SC* phase) or the chiral smectic H phase (hereinafter abbreviated to SH* phase) of the ferroelectric liquid crystal, and those having these phases in the vicinity of room temperature are preferred.\nThese chiral smectic liquid crystal materials may be obtained by blending a plurality of single compounds each exhibiting a chiral smectic phase by itself, but it is known that the materials may be also obtained by adding an optically active liquid crystal compound, preferably a chiral smectic liquid crystal compound to an achiral smectic liquid crystal (exhibiting smectic C phase (SC phase), smectic H phase (SH phase), etc.).\nVarious kinds of compounds exhibiting SC phase have been known, but as to whether or not chiral smectic liquid crystal materials obtained by adding an optically active liquid crystal to the above-mentioned known compounds exhibit superior performances in liquid crystal display utilizing ferroelectricity, no ultimate evaluation thereof has been yet obtained. This is due to the fact that liquid crystal display utilizing ferroelectricity has not yet been technically completed. Thus it is necessary in the present situation to test various novel SC materials."}
-{"text": "In recent years, optical devices, such as light-emitting diodes (LED), laser diodes, and UV photo-detectors have increasingly been used. Group-III nitride compounds, such as gallium nitride (GaN) and their related alloys, have been known suitable for the formation of the optical devices. The large bandgap and high electron saturation velocity of the group-III nitride compounds also make them excellent candidates for applications in high-temperatures and high-speed power electronics.\nDue to the high equilibrium pressure of nitrogen at typical growth temperatures, it is extremely difficult to obtain GaN bulk crystals. Therefore, GaN layers and the respective LEDs are often formed on other substrates that match the characteristics of GaN. Sapphire (Al2O3) is a commonly used substrate material. FIG. 1 illustrates a cross-sectional view of an LED package component. LED 2, which includes a plurality of GaN-based layers formed on sapphire substrate 4. Sapphire substrate 4 is further mounted on lead frame 6. Electrodes 8 and 10 electrically connect LED 2 onto lead frame 6 through gold wires 12.\nSapphire has a low thermal conductivity. Accordingly, the heat generated by LED 2 cannot be dissipated efficiently through sapphire substrate 4. Instead, the heat is mostly dissipated through the top end of LED 2 and through gold wires 12. However, the necessary length of gold wires 12 to extend to lead frame 6 renders the heat-dissipating efficiency low. In addition, electrode 10 occupies chip area and reduces the total chip area available for LED light output."}
-{"text": "As the display technology becomes increasingly mature, various displays have gradually developed. At present, liquid crystal displays (briefly known as LCDs) are more and more widely applied due to advantages such as little power consumption, miniaturization, light weight and thin thickness, and so on.\nAn existing liquid crystal comprises two polarizers having light transmitting axes perpendicular with each other and a liquid crystal panel located between the upper and lower polarizers. The principle of greyscale display is that the lower polarizer converts a natural light into linearly polarized light; the deflection status of liquid crystals is controlled by voltages so as to convert the linearly polarized light into elliptically polarized light, and the upper polarizer analyses polarization of the elliptically polarized light so as to realize different greyscale displays."}
-{"text": "Shoppers in retail stores typically use shopping carts or baskets owned by the store to hold their purchases as they travel through the store. Purchases are typically bagged at the check out counter and either reloaded into the cart or the bags themselves are carried out of the store by the shopper. Carrying the bags is cumbersome for the shopper and the number of bags that a shopper can carry is limited by the strength of his or her arms. With regard to using shopping carts outside of the store, such an arrangement does not work for city shoppers who walk to and from the store. In addition, even if the shopper drives to the store, and carries the groceries to his or her car with a store-owned cart, the store must go to the trouble of collecting the carts from the parking lot. In addition, the carts must be maintained, repaired and stored, which can add up to significant costs and take up valuable storage space.\nConsumer-owned folding wire carts for use in shopping are known. The shopper brings the cart to the store, opens/unfolds it, and uses it to hold items during shopping. The shopper then uses the same cart after checkout to transport the purchases to his or her car or home. The disadvantage of such carts, however, is that they can be heavy due to the metal wire construction. In addition, they tend to feature a very utilitarian appearance. They also typically do not fold completely flat, and thus take up valuable space which is undesirable (especially when the shopper is an apartment dweller).\nIn view of the above, a need exists for a collapsible personal trolley that is inexpensive, lightweight, durable and attractive."}
-{"text": "Today, people often utilize computing devices for a wide variety of purposes. Users can use their computing devices, for example, to communicate with other users. Such communications are increasingly popular over a social networking system. Digital communications, such as those on a social networking system, may involve various types of communication. Some types of digital communication allow a user to engage in focused exchanges. For example, the user may target a particular user or users through the use of a messaging system or an email system supported by a social networking system. As another example, the user can enter into audio communications or video communications with other users.\nIn many instances, video communications can be preferred by users because video communications can allow the users to most effectively convey information and simulate real life communications. In some instances, two participants in different locations can engage in video communications. It also can be desirable to allow a group of users in multiple locations to use video communications to facilitate communications among the group."}
-{"text": "This specification relates to online content provisioning.\nThe Internet provides access to a wide variety of resources, such as video and/or audio files, and web pages for particular subjects. Access to these resources has provided opportunities for content items to be provided with the resources. There are many different types of content items that can be provided, such as video items, text items, audio items, image items, and combinations thereof. For example, some online content items include a text portion and an image portion. The performance of an online content item is often affected by the form of media used to present the content item.\nSome advertisements only include a textual portion. In such cases the content item could be enhanced by accompanying the textual portion by an image portion to display jointly to users. However, unless a content item provider specifies an image or images to be provided with the textual portion, selecting a suitable image to display is a time-intensive process that involves a manual search of images and review of the images and the textual portion."}
-{"text": "1. Field of the Invention\nThis invention relates to circuitry for use with a guidance system for an automatically guided vehicle (AGV) and more specifically to circuitry for detecting the presence of a valid guidepath signal from various types of sensors, such as photocell arrays, infrared light, sound, and the like which operates over a large dynamic range without the need to make gain or threshold adjustments.\n2. Description of the Prior Art\nAutomatically guided vehicles are generally known in the art. Examples of such vehicles are disclosed in the following U.S. Pat. Nos.: 4,328,545; 4,345,662; 3,379,497; 4,623,032; 4,602,334; 4,627,511; 3,039,554; 3,610,363; 3,933,099; 4,003,445; 4,151,526; 4,602,334 and 4,500,970. Such automatically guided vehicles are used in a wide variety of applications, including the transfer of raw materials and subcomponent parts in manufacturing and assembly facilities, the cleaning of floors in warehouses and parking lots and the delivery of mail in business offices. AGVs are also used in a wide variety of applications in the agricultural industry such as, plowing, harvesting, mowing, and the like.\nIn each application the AGV is guided by on-board sensors which follow a guidepath. Various types of systems are used for automatically guiding the vehicle along a guidepath.\nIn U.S. Pat. No. 4,003,445, the guidepath consists of a fluorescent material, applied to a floor, carpet or the like which emits visible light in a predetermined frequency range. The fluorescent material, however, is normally invisible under ambient lighting conditions. In this type of guidance system, an ultraviolet light located on the vehicle, irradiates the guidepath. This radiation stimulates the fluorescent materials in the guidepath causing it to emit visible radiation which is sensed by sensors onboard the AGV.\nOptical guidance systems for use along a predetermined guidepath are susceptible to spurious operation due to background radiation. Thus, the guidance system must be able to detect a \"no-line\" situation. Also such systems are affected by the non-uniform intensity of the fluorescent guidepath. The solution to this problem has heretofore been attempted. For example, U.S. Pat. No. 4,003,445 discloses an automatic gain control (AGC) circuit for distinguishing between varying intensities of the reflected light from the fluorescent guidepath and background radiation. In this system, AGC feedback circuitry is disclosed for continuously adjusting the gain of the sensor circuit to compensate for the variations in the guidepath intensity. The AGC circuitry is also coupled to other circuitry for detecting a no-line situation. However, such circuitry is relatively complicated and requires a continuous gain adjustment of the sensor circuit to compensate for the varying intensity of the guidepath."}
-{"text": "A tensioner lever is generally mounted, by means of a mounting bolt or pin, on an engine block or other frame, in the vicinity of a tensioner which cooperates with the lever. The tensioner lever, in cooperation with the tensioner, maintains appropriate tension in the transmission medium to prevent transmission failure due to excess tension and excess loosening of the transmission medium.\nIn one well-known example of a conventional tensioner lever, described in Japanese Laid-open Patent Publication No. 2001-323976 (pages 1 to 4, FIG. 2), a sliding contact portion extending along the traveling plane of the chain, and a reinforcing body, which reinforces the sliding contact portion, are fused together. In another example of a conventional tensioner lever, described in Japanese Utility Model Registration Publication No. 2540896 (pages 1 to 3, FIG. 1), the lever includes a resin shoe for sliding contact with the transmission chain, and an aluminum arm for supporting the resin shoe.\nIn the first of the above two conventional chain levers, the reinforcing body is strengthened by glass fibers. As the reinforcing body continues to pivot about a pivoting shaft on an internal combustion engine, the internal surface of its mounting hole is subject to wear, and the glass fibers are exposed and crushed. The crushed glass fibers act as an abrasive, causing accelerated wear and damage to the mounting hole of the reinforcing body and the pivoting shaft. In the case of a lever incorporating an aluminum arm, the continued pivoting of the lever about a pivoting shaft on an engine causes the arm to be burned and thereby also subjected to wear and damage. The use of a bushing or the like, fitted to the pivoting hole, has been considered to avoid the problem of wear. However, this measure increases the number of parts, the difficulty of assembly, and the production cost. Moreover, in the case of an aluminum arm, recycling of a spent lever is troublesome because the resin shoe, the aluminum arm, and a resin pad, must be separated before disposal.\nAccordingly, the objects of the invention are to solve the above-described problems, and to provide a tensioner lever having excellent mechanical strength and wear resistance, reduced production cost, reduced weight, and ease of recycling. A particular object of the invention is to reduce wear on the inner circumferential surface and the boss portion of the pivoting hole by which the lever is mounted on the engine block."}
-{"text": "As optical radar apparatuses capable of setting a detection area in a widthwise direction of a road surface, there are apparatuses such as those disclosed in Japanese Utility Model Laid-Open Nos. 59-117980 and 59-117981.\nIn the apparatus disclosed in Japanese Utility Model Laid-Open No. 59-117980, shown in FIGS. 7A and 7B, light from a light source is condensed to some extent by a lens to obtain a first detection range .theta.T1, and the light source is suitably moved to obtain a second detection range .theta.T2 wider than the first detection range .theta.T1 in a widthwise direction of a road surface by changing the degree of condensing of the lens.\nIn the apparatus disclosed in Japanese Utility Model Laid-Open No. 59-117981, shown in FIG. 8, light from a light source traveling through a lens is diffused by a prism to increase an expansion angle in a road surface direction, thereby setting a detection range wide in a widthwise direction of a road surface.\nIn the apparatus disclosed in Japanese Utility Model Laid-Open No. 59-117980 among the above-described conventional apparatuses, shown in FIGS. 7A and 7B, however, light cannot be emitted simultaneously for the first detection range .theta.T1 and the second detection range .theta.T2, and the detection range is set one-sidedly, since the detection range is changed by moving the light source.\nThe apparatus disclosed in Japanese Utility Model Laid-Open No. 59-117981, shown in FIG. 8, entails the problem of a reduction in the maximum detection distance because light is uniformly diffused by the prism. Because of this problem, it is not possible to meet a demand heretofore made for widening the detection range in a widthwise direction of a road surface in a short-distance area without reducing the maximum detection distance when an optical radar apparatus is used for an inter-vehicle control or an obstacle detecting apparatus.\nThe present invention has been achieved in consideration of the above-described problems, and an object of the present invention is to provide an optical radar apparatus capable of increasing the expansion angle in a road surface direction in a short-distance range without setting the detection range one-sidedly and without sacrificing the maximum detection distance."}
-{"text": "This invention is directed to the providing of a transcutaneous nerve stimulator which is designed to be utilized in T.E.N.S. (Transcutaneous Electrical Nerve Stimulation) therapy. T.E.N.S. therapy is based on a non-invasive, non-narcotic concept of pain management which is non-addictive, is not subject to abuse, and does not interact with drugs. T.E.N.S. therapy has already proven to be an effective modality in treating the organic pain problems associated with the following conditions: chronic lumbar and cervical strains or sprains, degenerating disc disease, degenerative arthritic disease, neuropathies, neuralgias, post-lumbar laminectomy syndrome, post-thoracotomy syndrome, bursitis, postphlebitis syndrome, phantom limb syndrome, and tension and migraine headaches.\nEarly attempts to suppress organic pain and other neurophysical effects utilizing electrical stimulation occurred as early as about 2,000 years ago when it was discovered that gout apparently cuold be successfully treated by placing the diseased extremities in a tub of water filled with electric eels. Later, headaches were treated using a similar approach. A detailed, scientific investigation was finally conducted by Professor Galvani of the University of Bologna, which investigation is credited with ultimately leading to the development in the 1800's of electrical equipment for suppression of organic pain.\nThe earliest therapeutic devices utilizing electrical simulation for the most part featured a constant amplitude and rate. Examples of the early art are Benz, U.S. Pat. No. 646,793; Raymond et. al., U.S. Pat. No. 872,148; Tibbals, U.S. Pat. No. 1,059,090; and Call, U.S. Pat. No. 1,908,688. A major problem with electrical stimulation therapy was, and still is, accommodation, whereby the nerve being stimulated in effect accommodates itself over time to the electrical charge, such that the effectiveness of the treatment is diminished.\nIt took scientists a long time to discover, and attempt to address the problem. Nemec, U.S. Pat. No. 2,622,601; DiPerma, U.S. Pat. No. 2,624,342; and Gratzl, U.S. Pat. No. 2,771,554 all disclose electrotherapeutic devices with at least one including means to vary the rate, amplitude or pulse width of the generated electrical pulse. However, merely being able to change either the rate, amplitude or pulse width still resulted in the problem of accommodation occurring, unless an individual manually adjusted the controls prior to the occurrence of accommodation. The process was both labor intensive and inefficient, with respect to the quality of the therapy, since maximum pain relief was not being provided.\nIn 1967, a Dr. Sweet at Massachusetts General Hospital developed the first T.E.N.S. unit. The effectiveness of T.E.N.S. therapy is based on its incorporating two major pain control theories. Under the so-called Gate Control Theory, pain can be inhibited and suppressed by \"closing the gate\" on pain signals as such signals arrive at Central Nervous System centers. This theory postulates that by providing electrical stimulation of a sufficiently high amplitude, the electrical signals race up large myelinated fibers faster than the pain signals travel up smaller myelinated or unmyelinated fibers. The neutral impulses transmitting pain information to the brain thus become interrupted, and since the brain fails to receive the pain signals, no pain is perceived. The other theory incorporated in T.E.N.S. units is the Endorphin Theory, also known as the Endogenous Opiate Theory. This theory postulates that the sustained input of T.E.N.S. signals triggers the release of naturally occurring pain making endorphins and enkephalins (morphine-like substances). These natural substances seemingly block pain signals by a mechanism similar to conventional drug therapy, and inhibit pain information from reaching the brain.\nHowever, T.E.N.S. units, like all electrotherapeutic devices, have suffered problems with accommodation. For example, Geerling, U.S. Pat. No. 4,019,519 issued in 1977, disclosed a unit having only its amplitude adjustable. Miller, U.S. Pat. No. 4,084,595 issued in 1978, disclosed a unit having its rate, amplitude, and pulse width all independently, manually adjustable. However, even this advancement resulted in a less than efficient treatment of the problem of accommodation, since either the therapist or the patient had to, in theory, repeatedly adjust the controls in an attempt to avoid accommodation.\nAlthough variation enables one to deal with the problem of accommodation, pain relief is sacrificed. This is in part due to the interaction between amplitude and pulse width. There is a clinical correlation between amplitude and pulse width with regard to the efficacy of the stimulus. As one shortens the duration of a pulse, its amplitude must be increased to maintain the efficacy of the stimulus. This relationship when plotted graphically is known as a strength-duration curve. Thus not only must the ideal T.E.N.S. units have adjustable amplitude and pulse width, but it must also be able to modulate those values in such a way as to approximate the A-gamma-motor strength-duration curve.\nIn addition to amplitude and pulse width approximating the strenght-duration curve, the rate of the pulses must vary so as to eliminate any potential for accommodation. This explains the short-comings in Reiner, U.S. Pat. No. 2,808,826 which disclosed a unit which permitted instantaneous changes in pulse width and amplitude to two pre-set points along the strength-duration curve, and Maurer, U.S. Pat. No. 4,340,063 which disclosed a unit having its amplitude modulate in response to modulations in pulse width so as to approximate a portion of the strength-duration curve. The rate in Maurer was adjustable, but only to the extent taught by Miller, such that the problem with accommodation still existed."}
-{"text": "To satisfy the increasing demand of bandwidth and speed, pluggable transceiver modules are used in line cards on various network devices (e.g., switches, routers, etc.). The pluggable transceiver modules are used to convert electrical signals to optical signals or in general as the interface to a network element copper wire or optical fiber. There are tests that need to be performed on the interface between the transceiver and a host for compliance with standards and to ensure that the transceiver will work properly with the host. Conventional testing is time intensive and involves high costs, and may not be scalable in a manufacturing and production environment.\nCorresponding reference characters indicate corresponding parts throughout the several views of the drawings."}
-{"text": "1. Field of the Invention\nThe present invention relates to an adaptive writing method for a high-density optical recording apparatus and a circuit thereof, and more particularly, to an adaptive writing method for optimizing light power of a light source, e.g., a laser diode, to be suitable to characteristics of a recording apparatus, and a circuit thereof.\n2. Description of the Related Art\nWith the multi-media era requiring high-capacity recording media, optical recording systems employing high-capacity recording media, such as a magnetic optical disc drive (MODD) or a digital versatile disc random access memory (DVD-RAM) drive, have been widely used.\nAs the recoding density increases, such optical recording systems require optimal and high-precision states. In general, with an increase in recording density, temporal fluctuation (to be referred to as jitter, hereinafter) in a data domain increases. Thus, in order to attain high-density recording, it is very important to minimize the jitter.\nConventionally, a write pulse is formed as specified in the DVD-RAM format book shown in FIG. 1B, with respect to input NRZI (Non-Return to Zero Inversion) data having marks of 3T, 5T and 11T (T being the channel clock duration), as shown in FIG. 1A. Here, the NRZI data is divided into mark and space. The spaces are in an erase power level for overwriting. The waveform of a write pulse for marks equal to or longer than 3T mark, that is, 3T, 4T, . . . 11T and 14T is comprised of a first pulse, a last pulse and a multi-pulse train. Here, only the number of pulses in the multi-pulse train is varied depending on the magnitude of a mark.\nIn other words, the waveform of the write pulse is comprised of a combination of read power (FIG. 1C), peak power or write power (FIG. 1D) and bias power or erase power (FIG. 1E). Here, the respective power signals shown in FIGS. 1C, 1D and 1E are all low-active signals.\nThe waveform of the write pulse is the same as that in accordance with the first generation 2.6 GB DVD-RAM standard. In other words, in accordance with the 2.6 GB DVD-RAM standard, the waveform of the write pulse is comprised of a first pulse, a multi-pulse train and a last pulse. Although the rising edge of the first pulse or the falling edge of the last pulse can be read from a lead-in area to be used, adaptive writing is not possible since the write pulse is fixed to be constant.\nTherefore, when a write operation is performed by forming such a write pulse as shown in FIG. 1B, severe thermal interference may occur back and forth with respect to a mark in accordance with input NRZI data. In other words, when a mark is long and a space is short or vice versa, jitter is most severe. This is a major cause of lowered system performance. Also, this does not make it possible for the system to be applied to high-density DVD-RAMs, e.g., second generation 4.7 GB DVD-RAMs."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to a nail gun, and more particularly to a safety structure of a nail gun.\n2. Description of the Related Art\nTypically, to prevent a nail gun from being triggered unexpectedly to hurt somebody, the nail guns are equipped with a safety for protection of user and people around."}
-{"text": "The field of art to which this invention pertains is the treatment of aqueous streams containing water-soluble inorganic sulfide compounds. More specifically, the invention is directed toward a method for treating an aqueous stream containing a water-soluble, inorganic sulfide compound."}
-{"text": "In general, conventional rate converters include two major computation blocks, as illustrated in FIG. 1. A first block is a conversion rate tracking loop that determines the ratio between an input sample rate and output sample rate. Once the ratio is determined, the conversion rate tracking loop generates the corresponding input sample index for each output sample.\nThe second computation block is the sample interpolator. The function of this block is to interpolate the input data sample and to generate the output data sample with the real value index. One problem with conventional rate converters is that they create aliasing from the resampling, and it is relatively difficult to produce an output signal having a high Signal-to-Noise Ratio (SNR) and having a relatively flat frequency response within standard frequency ranges for audio signals.\nEmbodiments of the invention address this and other limitations of the prior art."}
-{"text": "Refrigerator appliances generally include a cabinet that defines a chilled chamber for receipt of food items for storage. For example, the cabinet can define a fresh food chamber and a freezer chamber. The fresh food chamber can be maintained at a temperature greater than the freezing point of water. Conversely, the freezer chamber can be maintained at a temperature equal to or less than the freezing point of water.\nRefrigerator appliances generally also include one or more shelf assemblies positioned within the fresh food chamber and/or the freezer chamber to facilitate storage and/or organization of any food items positioned therein. Typically, the shelf assemblies are adjustable in height using, e.g., a cantilevered track assembly. Accordingly, with such a configuration, the user may customize the fresh food chamber and/or the freezer chamber of the refrigerator appliance to meet their specific needs.\nModern refrigerator appliances can also include lighting systems attached to or embedded within the one or more shelf assemblies. The lighting systems can, e.g., illuminate the shelf assembly itself, or alternatively can illuminate a lower shelf assembly. However, in order for such refrigerator appliances to provide such lighting systems with electrical power, one or more electric wires are generally required to be disconnected and reconnected as the shelf assembly is adjusted. Certain consumers may find such electrical wires unsightly and may find the additional steps of disconnecting and reconnecting wires undesirable.\nAccordingly, a refrigerator appliance having one or more shelf assemblies capable of connecting a lighting system to an electrical source without requiring disconnection and reconnection of electrical wires would be useful. Moreover, a refrigerator appliance having one or more shelf assemblies capable of connecting a lighting system to an electrical source without any visible connections would be particularly beneficial."}
-{"text": "1. Field of the Invention\nThe present invention relates to the filed of the digital signal processing, and more specifically, to the field of phase equivocation in burst modems.\n2. Discussion of the Prior Art\nWhen the communication exists for a long period of time, it is possible to use a conventional phase-lock loop (PLL) in a receiver to recover the carrier from a continuously transmitted signal. The PLL that is specifically designed to recover the carrier is termed as a phase carrier loop.\nHowever, in a time-division multiple access (TDMA) communication system, the available spectrum is used by letting each user to have access to the whole band for a short time (traffic burst), during which time the user transmits data as fast as the user can. The user's frequency allocation is shared with the other users who have time slots allocated at other times.\nThus, in a TDMA-type burst communication systems, the signals exist for a short period of time. Consequently, there is little time in which the phase carrier loop can slew to the proper phase in order to recover the carrier from the burst signal.\nThis situation is exacerbated for a TDMA-type burst commination system if the phase of the incoming symbol is very close to the border between the plant points. If this is the case, the slicer in the phase carrier loop can not determine to which quadrant the incoming symbol point having such phase should belong because such symbol can be equally located in either of two quadrants. Such symbol point is further referred to as a symbol point having a phase equivocation, and the phase carrier loop can not lock fast enough on an incoming signal having a phase equivocation characteristic.\nWhat is needed is to design a carrier phase loop system that is capable of resolving the phase equivocation of the incoming symbols, thus placing any incoming symbol point to a certain quadrant for further processing by the loop, thus increasing the speed of locking of the carrier phase loop on the carrier."}
-{"text": "A variety of batteries employing metallic magnesium anodes are known. Morehouse et. al., in U.S. Pat. No. 2,759,986, taught a magnesium anode cup container lined with a cathode material paste consisting of powdered sulfur, graphite, and magnesium bromide electrolyte. Roche et. al. in Argonne National Lab, ERDA Report No. 76-35, \"Alternative Secondary Cell Systems\", Progress Report, 1976, pp. 53 to 54, discloses a molten electrolyte, magnesium-iron sulfide cell, using FeS.sub.2 cathodes, and MgCl.sub.2 -- NaCl -- KCl electrolyte. The cathode was covered with a stainless steel screen and the anode was covered with graphite fabric. The cell operated at 450.degree. C but short circuited after four cycles due to magnesium dendrite formation.\nOthers in the art have substituted iron sulfide as the cathode material in combination with a lithium anode. These lithium cells provide outstanding high voltage characteristics. Vissers et. al., in U.S. Pat. No. 3,933,521, taught a molten electrolyte, rechargeable, electrochemical cell which required a LiCl -- KCl electrolyte between the anode and cathode and operated at temperatures of about 400.degree. C. The cathode contains iron, cobalt or copper sulfides in a conductive support structure, insulated with a separate boron nitride or yttria cloth separator, and disposed in the center of a housing container. The anode which lines the bottom of the container comprises a 90% porous, compacted, stainless steel or nichel felt, coated with cobalt and impregnated with molten lithium.\nThe use of lithium, due to its high reactivity, requires the use of expensive separator and container materials. In lithium cells, the anode and electrolyte are very expensive. The use of magnesium metal provides severe short circuit complications due to dendritic growth at the magnesium electrode during recharging. Magnesium also reacts chemically with a wide variety of separator materials, particularly oxides, to make them electrically conducting, thus causing shorting of electrodes. What is needed, is an inexpensive, secondary, electrochemical cell, which will still provide good voltage characteristics, and which will not be easily degraded over long time periods."}
-{"text": "The present invention relates to a multiple pocket, expandable envelope, a blank for forming a multiple pocket, expandable envelope, and a method of forming a multiple pocket, expandable envelope from a blank. More particularly, the invention relates to the use of a planar, unitary blank for forming a multiple pocket, expandable envelope which may be folded simply and inexpensively in a one pass, right angle gluing and folding operation.\nIn constructing envelopes or folders, it is often necessary to provide a plurality of individual pockets or compartments for separating the contents of the envelope. Additionally, the envelope must be capable of being shipped and stored in a flat, collapsed configuration to use shipping and storage space efficiently, and then expanded to receive varying amounts of contents, e.g., sheets of paper.\nConventional multiple pocket envelopes or folders are formed from a plurality of components which are automatically and manually assembled. The forming of such envelopes from a plurality of components is difficult and expensive, involving both manual steps and complex machinery.\nTo form multiple pocket, expandable envelopes simply, inexpensively and efficiently, such envelopes must be formed from a planar, unitary blank which may be folded and glued in one pass, right angle gluing and folding operation. Such operation eliminates manual handling and permits the envelope to be manufactured on relatively simple, conventional folding and gluing apparatus."}
-{"text": "A digital transport stream typically comprises scrambled audio/video/data content packets and scrambled conditional access messages (ECM, EMM) which have to be processed in order to extract control data (CW); the latter being necessary for the multimedia unit to descramble the content packets relating to the event (e.g a broadcasted program or a set of programs). Conditional access messages include two kinds of information, namely Entitlement Control Messages (ECM) and Entitlement Management Messages (EMM). The ECM is information relating to content packets, whereas the EMM is information dedicated to an individual end user (subscriber) or a group of end users. The ECM contains the access conditions for the current broadcast event together with control word (CW) for unscrambling this event. The Control Word (CW) is a key used for decrypting content packets of a packetized event broadcasted within the data stream. Thus, each ECM is specific to each event (e.g. a TV program). The EMM is a message used typically by the security module to set, reset or change product access entitlements, credit, etc. . . . Thus EMM refers to the rights (access data) of the subscriber for accessing to the content (events).\nThere is two main types of set-top-boxes known by the prior art. The first type comprises set-top-boxes that can generally receive broadcast unencrypted events, and which therefore do not require any access control. As internet access is widespread at user homes, such set-top-boxes are increasingly provided with an IP-connection which allows the end user to use his set-top-box together with his television screen as a terminal for internet facilities.\nThe second type refers to set-top-boxes provided each with a security module. These set-top-boxes are typically used by Pay-TV systems for processing selective access to broadcast services on a per-subscriber basis managed by a conditional access system (CAS). Such a security module typically refers to a smart card (chip card) which can be inserted into a slot of a conditional access reader. This card reader can be connected to the multimedia unit or be lodged directly within this unit. Thus, the security module becomes a necessary device for decrypting the scrambled content packets of the digital transport stream. To this end, the conditional access messages contained in the transport stream are routed to the security module which has the keys to decrypt the ECM in order to extract the control word (CW) which then will be used by a descrambling module (within the multimedia unit) for descrambling the content packets. This process is achieved only if the conditional access rights (provided by the EMM) are valid and checked by the security module.\nThe problem with set-top-boxes of the first type is that they are often unsuitable for processing conditional access messages required for descrambling audio/video content packets within a Pay-TV environment. This is true even if they are provided with an IP-connection. Thus, a user wanting to subscribe a subscription for having access to Pay-TV programs has to change his set-top-box with another one that supports conditional access processing.\nOn the other hand, the set-top-boxes of the second type, provided with a security module, are dependent on this module for descrambling the content packets of the transport stream. It means that the security module of the set-top-box is required for descrambling purposes. If this security module becomes unavailable (for instance because the required smart card has to be changed to a renewal or for any other reasons), this set-top-box becomes at least temporarily unavailable for receiving Pay-TV programs.\nAccordingly, there exists a need for improving processes and multimedia units (e.g. set-top-boxes) in order to optimize the processing of control data for descrambling the content packets of the transport stream, for instance by adapting this processing to different circumstances, depending on the availability and the adequacy of the network and/or depending on the material configuration at the end user."}
-{"text": "Drums are the simplest, and most primitive, musical instruments. The drum is a percussive instrument, which produces sound by striking a membrane. The sound is propagated through a membrane, or drumhead, to the drum shell, which is designed to resonate when the membrane is struck. The drumhead is coupled to the drum shell through drum hoops, lugs, and lug or hoop holders. The energy created by striking the drum head is transferred into a wave in the drum shell, producing the distinctive drum sound, a tuned impulse.\nWhile drums usually cannot play different pitches, they are usually tuned. The drum is tuned by tightening or loosening the drumhead by adjusting the lugs and drum hoops. The tighter the drumhead, the higher the pitch propagated by the drumhead.\nMany drummers use a drum kit. Drum kits have several different drums, which can be individually tuned. A drum kit is often composed of various drums, such as a bass or kick drum, snare drums, and tom drums, as well as assorted cymbals and high-hats. When a drummer is drumming, there is substantial vibration throughout the drum kit. Additionally, the various drums can move or flex as they are struck, meaning that the drums, themselves, are vibrating and, therefore, moving. This is especially true of the tom and snare drums.\nA drum mount or support is a sub-classification that includes many different methods of mounting the drum to legs or to other structural elements. The mounts used for drums often degrade the sound, because the drum is held too tightly, damping or attenuating the tuned impulse. Depending on how the drum mount is attached, and how it supports the drum, it can hinder the drum shell resonance, the drumhead, or both. Any drum mount or support that rigidly fixes itself to the drum, whether to the drum shell or the drum hoops, risks damping the sound. A very small number of mounts attach to lugs or lug holders. As currently implemented, these integrated mounting brackets and lugs are a bad idea, because they fail to properly isolate the tuning of the lug from the tuning of the drum mount. The lugs are supposed to be tuning the drum. By adding additional stiffness (force) to the lugs, a fixed drum will change the tuning of the drum. An integrated lug and drum mount is conceptually attractive, because such a device, if it did not adversely affect the tuning of the drum head, would be easier to install on the drum.\nOther mounts attach to the top and/or bottom hoops of the drum shell. Still other supports attach to the drum shell. Rigidly attached supports, whether mounted to the shell or hoops, will damp the vibration of the drum, and may, ultimately, distort the sound through buzzing or rattling, if the support mount is not properly engineered and attached.\nAs a result, a new device is needed, integrating the lug and drum mount, without unduly affecting the tuning of the drum. The new integrated lug and mount should allow the drum shell to resonate freely. Such an integrated lug and mount needs to maintain the drum, in position, without, itself, generating objectionable sounds. The integrated lug and drum mount should be quickly and easily adjustable. It should also allow for quick set-up and break-down of a drum kit. An integrated lug and drum mount that has variable stiffness, and isolates the drum shell from the structural members supporting the drum would be ideal."}
-{"text": "Wavelength-division-multiplexing (WDM) is a key enabling technology in today's high speed digital communication infrastructure that supports vast amount of data transportation which in turn is essential for many data centric informational applications such as, for example, many internet based applications. The transportation of vast amount of data are made via optical digital signals over an extensive optical fiber network across the country or global and such optical signals, riding on different wavelengths (or \u201ccolors\u201d), are distributed and/or re-distributed among various parts or branches of the optical network via signal exchanges located at various data centers and other facilities. Sometimes, such signal distribution and/or re-distribution may also involve conversion of signals from optical to electric, and then from electric to optical, in terms of its carrying media.\nWhen optical signals are distributed and/or re-distributed within an optical network, it often involves interconnecting optical signals from one signal handling unit, which may be, for example, an optical signal transponder installed in a shelf hosted by a rack (\u201cbay\u201d), to another signal handling unit which may be located in a same room, in a different room of a same floor, or sometimes in a different floor. FIG. 15 is a demonstrative illustration of an interconnected optical system 900 as is known in the art. System 900 is simplistically illustrated to include a first bay 910 and a second bay 920 located, for example, in a same room of a building. Bay 910 may include, from top to bottom, multiple shelfs such as shelf 911, with each of the shelfs having multiple optical signal transponders.\nAn optical signal coming from an optical signal transponder located in bay 910 may be connected to another optical signal transponder located in bay 920 via a piece of fiber 902, which may have connectors 901 and 903 at its two ends connecting to the signal transponders. Multiple optical signals from bay 910 may need to be connected to multiple destinations in bay 920, and vise versus, using multiple pieces of fibers. Generally the number of fibers needed equals to the number of optical signals being interconnected between the bays, which is demonstratively illustrated in FIG. 15 by a second piece of fiber 904 and the \u201cdots\u201d in between which represents the existence of many more fibers between bay 910 and bay 920.\nWith the ever increasing data rate, in particular rapid deployment of WDM technology, the number of optical signals of different wavelengths that need to be interconnected between different bays, and sometimes between different shelfs in a same bay, has increased dramatically resulting in the explosive use of fibers in signal interconnection. FIG. 16 is an exemplary picture of a traditional optical system 990 where optical signal interconnect among the bays are provided by an enormous stack of individual pieces of fibers 991 which, like cooked spaghetti, are usually \u201cdumped\u201d at the backside of the various bays."}
-{"text": "This invention relates to a method for manufacturing a plywood using an improved veneer bonding technique.\nIn prior art techniques of bonding veneers for plywood manufacture, it is well known that poor adhesion such as puncture and deficient bond strength largely depends upon the extent of drying of veneers, that is, the moisture content of veneers.\nHowever, it is impossible at the present state of the art to dry a number of veneers to an equal moisture content, and veneers are dried to a more or less varying extent. Particularly, the recent deterioration of log stock results in a mixed supply of difficult- and easy-to-dry woods. With the increased extent of drying, over-drying results in excessive shrinkage and corrugation of veneers. With the reduced extent of drying, on the other hand, under-drying results in poor adhesion. In connection with veneer drying, a number of problems arise at the site of manufacture.\nThe only one solution for the present status depends on the development of a practical gluing technique which is completely or substantially independent from variable moisture contents. To essentially eliminate these problems, it is necessary to establish a technique capable of gluing veneers having a high moisture content at low cost.\nOne technique of gluing veneers having a high moisture content is disclosed in Japanese Patent Application Publication No. 54-3929 titled \"Plywood Manufacturing Method\". According to the method of this patent publication, a plywood adhesive fluid of well-known formulation is applied to a veneer and dried thereon, and another veneer is placed on the adhesive-applied veneer followed by hot pressing. Although the results described in the patent publication are satifactory, I found that this technique has the following problems. With a dryer equipped for the purpose of drying the adhesive applied and set to proper drying conditions including temperature and time, the adhesive applied on those veneers having a lower moisture content is overdried and the adhesive applied on those veneers having a higher moisture content is under-dried, failing to dry the adhesive applied to a desired extent for all veneers. As a result of the varying dryness of the adhesive, the bond strength achieved by hot pressing also varies over a wide range.\nIn the examples described in the patent publication, drying conditions are changed in accordance with the moisture content of veneers to be handled, and more specifically, the higher the moisture content of veneers, the higher the drying conditions are set. The dryness of the adhesive applied not only depends on the selected adhesive drying conditions, but largely depends on the moisture content of veneers. The varying bond strength is attributable to this fact. Consequently, the dryer for drying the adhesive applied to veneers must be of size and capacity as large as the conventional veneer dryers. Besides, it is well known that there is a significant difference between sapwood and heartwood portions of the same log and between spaced portions of the same veneer. From a standpoint of view of manufacturing plywoods of consistent quality from a variety of starting logs, the above-mentioned shortcomings of the prior art techniques are difficult to overcome in the actual forms of plywood manufacture requiring large-scale production and prevent commercial application of such techniques.\nAn object of this invention is to eliminate the above-mentioned shortcomings involved in the prior art techniques and to provide a novel and improved technique for firmly bonding veneers, even highly wet veneers if necessary, at low cost independent of the moisture content of veneers so that the technique may be readily applied to commercial manufacture."}
-{"text": "Solar energy appears to be the single most attractive source of energy. Solar systems are not liable to wind damage, as is the case with wind energy systems except in extreme weather conditions, which are exceptional, and would damage even conventional energy transmission systems. Solar collector can be installed almost anywhere, unlike tidal energy systems which must be located in coastal areas. Solar energy systems consume no fuel and the energy is, therefore, \"free\", in a sense. It causes no emissions to pollute the environment. It does not appear to require complex transmission and delivery systems such as electrical power lines, gas or oil pipe lines, or other delivery systems, provided it can be collected \"on site\", i.e., where it is to be used.\nSolar energy has however one major disadvantage, namely, that it is available only intermittently, in the hours of daylight. Accordingly, it is obvious that sufficient solar energy must be collected during the hours of daylight, that a sufficient quantity of energy can be stored, for reuse during the hours of darkness when solar energy is not available. Clearly, systems could be engineered to carry out these functions on a very large scale. However, that would then require the collection and storage of large quantities of solar energy which would then have to be transmitted from place to place, for example, to a large number of homes, connected to a central solar collector. This would result in excessive capital costs, and would likely result in relatively high energy transmission losses. Clearly, such centralized systems would fail to take advantage of the principal significant fact concerning solar energy, namely, in that it is substantially universally available over a large portion of the face of the earth, and thus does not require special transmission systems, to transmit the energy from one place to another. The theoretically ideal solar energy system would be self-contained, so that individual systems could be provided for individual buildings or homes, so that each building or home could be essentially energy self-sufficient. In this way, complex transmission or storage systems would not be required.\nIt is however apparent that in the design of any solar energy collector system it will be necessary to collect solar energy falling on a relatively substantial surface area, to provide sufficient energy even for a single home or building. In order to achieve this, numerous proposals have been made for providing solar collectors based on some form of concave curved mirror usually of parabolic shape having a relatively large surface area. The mirror catches the sun's rays falling on all points of the mirror, and concentrates them into a single point, resulting in a high energy concentration of light at that point. Usually, at the focal point of the mirror, a secondary reflector is provided for reflecting the concentrated point of light, and redirecting it, so as to convert the energy of the concentrated light beam into some other form of energy. Typical proposals of this general type are shown in, for example, U.S. Pat. Nos. 4,286,581, 4,295,462, 4,340,031, and 4,696,285. It is, of course, implicit in any of these solar collector devices that the actual mirror itself must be located in the open air. In all of these proposals, the mirror concentrates the solar energy falling on the mirror at the focal point of the mirror. The concentrated light is re-directed into a portion of the collector where the energy of the light beam is then converted into heat. For example, in U.S. Pat. No. 4,286,581, the solar energy is utilized to heat a tubular chamber, or cavity, and a fluid medium is passed through the chamber or cavity, and absorbs the heat from the concentrated light rays. The energy in the fluid is then transported to some other location for use. In these proposals the mirror was usually mounted on some form of a suitable equatorial tracking mount, so that the mirror could always be pointed directly at the sun. Accordingly, in the proposals of the type described above, it was usually necessary for the entire apparatus consisting of the mirror, the secondary reflector, and also the heat conversion chamber, all to be mounted together and moved simultaneously to track the path of the sun. Such systems therefore require highly complex suspension and tracking mechanisms, to provide the facility for accurate tracking of the sun, by all of the components simultaneously and also require heat transmission systems for transmitting the heat from the outside to the inside of a building.\nOther proposals for collecting solar energy have been made of a more simplistic nature, for example as shown in U.S. Pat. No. 4,059,226. In this type of proposal, large chambers are filled with any suitable heat storage medium. In this case, the medium is simply rock or stone. The chambers are fixed in position, and have glass panels through which the sun's rays may pass and heat up the stones. In the particular example, of this patent, insulating panels are adapted to be swung down over the glass panels, when the heat generated by the sun's rays falls below a certain preset minimum temperature.\nSystems of this kind are relatively massive installations, and are suitable in only relatively limited \"sun belt\" portions of the globe, and must obviously be located for optimum operation in a fixed location, i.e., facing south in the northern hemisphere.\nAnother proposal is shown in U.S. Pat. No. 3,988,166. In this patent, a concave mirror is provided, mounted on a suitable tracking mount, and having a secondary light concentration mirror at the focal point. The light falling from the main mirror on the secondary mirror is concentrated into a light beam. A central opening in the main mirror allows this concentrated light beam to pass through the main mirror. An energy collection chamber is provided immediately behind the main mirror, for receiving energy from the concentrated light beam. The converted energy is then transmitted from this chamber, to a location where it may be of use. While this system is somewhat simpler to construct, it still has the disadvantage that the energy conversion takes place out of doors, and requires the transmission of an energy conversion medium, typically, for example, a fluid medium, from the solar collector itself into the building or facility where the energy is required, typically to heat the building or to provide hot water for a domestic hot water system. Thus thick insulation must be provided to limit heat losses in transmission.\nIn a further proposal, in U.S. Pat. 4,373,514 granted on Feb. 15, 1983 to Lois, there is shown a form of solar collector system in which the collimate beam is routed through a conduit, around various angles. This system uses at least three separate mirrors, to bring the light beam along the angled conduit, to a point where it could be converted. However Lois does not disclose an equatorial mount for tracking the sun's path each day, and each day of the year, and consequently would have very little if any efficiency. In addition, the use of three separate mirrors within the conduits produces further heat transmission losses, which accumulate from mirror to mirror. Consequently, the system would be impractical for the purposes of for example domestic heating or heating of individual buildings, and would be extremely expensive."}
-{"text": "There have been an increasing number of machines activated by cordless power supply utilizing electromagnetic inductance, such as electric toothbrushes, cordless telephones, and portable devices (e.g., PTL 1). There have also been developments of machines activated by cordless power supply utilizing magnetic resonance, in relation to wall-hang television sets and personal computers (e.g., PTL 2). To add this, there have been many developments and suggestions for magnetic elements for wireless power transmission capable of feeding high power with high power transmission efficiency, in the field of wireless power transmission.\nHowever, feeding high power with high power transmission efficiency in such a magnetic element for wireless power transmission causes an excessive heat generation. This generated heat in the magnetic element for wireless power transmission may cause a problem to the magnetic element itself or to the functional parts of the element.\nTo address this problem of heat generation, for example, one approach is to adopt a structure of substrate having an in-built coil such as the one disclosed in PTL 3, in which a planar coil conductor to serve as the heat source is covered with a magnetic layer so as to radiate the heat outside through a conductor layer for heat transfer which is provided to the magnetic layer. This, without a doubt, enables radiation of heat generated by the planar coil to the outside."}
-{"text": "Coupling devices have been used extensively for many years in the railroad industry to connect adjacent ends of a pair of railway cars together to form a train consist of several cars. On those railway cars which will be used in interchange service, however, such coupling devices must receive approval from the Association of American Railroads (AAR) before they can be installed on the cars. It has come to be quite well known, in this railway application, that such coupling devices perform a number of important functions. Obviously, an important function of the standard coupling device is that it facilitates the connection and disconnection of individual cars to and from, respectively, a train consist. Another important function of a standard coupler is that it enables such railway cars to negotiate the various curved portions of the track structure which are encountered during operation. Further, these coupling devices will enable a number of cars to be readily combined so that a train consist can be made up, or such cars can be easily separated individually for either loading or unloading cargo thereto or therefrom, respectively, as required. This will allow a railroad to leave a particular car at a customer's plant while delivering or picking up other cars at other locations.\nIn more recent times, the railroad industry has come to recognize, however, that a number of rather significant advantages can be achieved by the interconnection of several railway cars together to form a generally semi-permanent unit. This is particularly the case, for example, where such individual railway cars are adapted for use in what is commonly known in the railroad industry as \"piggyback\" service. A primary reason for this is that the cargo to be either loaded or unloaded is brought to or removed from, respectively, a central location which is owned and operated by the railroad normally. Generally, this cargo is either over-the-road trailers or very large containers to be shipped by sea. The individual cars which have been connected together in such substantially semi-permanent fashion are commonly known in the railroad industry as a \"10-pack\". These 10-pack units do not require the use of standard couplers except at each outer end of each 10-pack unit. The primary reason why such standard coupling devices are not required is that in view of their dedicated service these 10-pack units will only be broken periodically. For example, this will normally only occur when maintenance must be carried out on an individual coupler component or on some other component positioned on such railway cars that will necessitate such car be taken out of service at least temporarily. It has become obvious to the railroad industry that with the use of this coupling arrangement it is possible for them to achieve a significant reduction in their costs. Such cost-savings can be attributed to a number of reasons. For example, these reasons include lower equipment weight which results in enhanced energy savings and fewer railway trucks which results in both lower equipment cost as well as a significant reduction in maintenance requirements. However, now with the rather extensive use of these semi-permanent coupling arrangements, particularly with new cars being built for piggyback service, and with the ever increasing loads that are now being carried by modern railway cars and trains, it has been determined that it is of the utmost importance for a close-buttoned relationship to be maintained between the draft components of the coupler devices. Such close-buttoned relationship is necessary, for example, so that the effects of the impact forces, which are usually encountered under normal buff conditions during train operation, can be reduced to an acceptable level. In this manner, unnecessary damage to both cargo and rolling equipment can be held to a minimum.\nTaught in U.S. Pat. No. 4,258,628, is one prior art type of articulated coupling apparatus which can be used for the above-identified purpose of connecting adjacent ends of a pair of railway cars together in such semi-permanent manner. This particular articulated coupling arrangement, as taught and illustrated therein, includes a male connection member secured to one end of a first railway car body and a female connection member which is secured to an adjacent end of a second railway car body. However, as can be seen the ends of both the male connection member and the female connection member which are secured to such car bodies do not provide any flexibility in how they are attached to such car body. The female connection member is rotatably-engaged in a center plate bowl portion of the bolster of a railway car truck in this prior art arrangement. This rotatable engagement is carried out in a manner that is well known in the railway art. The outer end portion of the male connection member is disposed for movement within a cavity formed in the outer end portion of such female connection member. To connect both the male connection member and the female connection member together in such semi-permanent manner, a pin member is utilized. This pin member is positioned in a vertical direction and is disposed in axially aligned apertures which are formed in each of the male connection member and the female connection member.\nAs taught in this reference, the aperture which is formed in such male connection member for receiving the pin member therein must be somewhat larger than the pin member itself. This larger aperture is required so that, while in operation, certain required movements of the coupling arrangement can be achieved. Additionally, a rear surface portion of such aperture that is located in such male connection member and which receives the pin member therein has a horizontally disposed concave configuration and a vertically disposed convex configuration. This particular configuration is required because it enables both the male connection member and the female connection member to move in each of a horizontal direction and a vertical direction in relationship to one another. Furthermore, this configuration provides, at the same time, a relatively substantial area of surface contact between the rear surface area of such pin aperture and the pin member itself.\nThe outer end surface of the outer end portion of such male connection member is provided with a convex configuration. This convex configuration abuttingly engages a complimentary concave surface which is formed on a front face portion of a follower member. In this coupling arrangement, such follower member is carried within the rear portion of such cavity disposed in the outer end portion of such female connection member. A pair of vertically disposed, slot-like cavities are provided on such follower members adjacent the rear surface thereof.\nA first portion of a resilient member is disposed within each of such vertical slot-like cavities. A second portion of each such resilient member extends outwardly from such rear face portion of the follower member. In this manner, a portion of the exposed outer surface of each such resilient member can be engaged by a vertically disposed wedge-like element. This wedge-like element must be provided in this prior art coupling arrangement to urge both the follower member and the male connection member forward.\nWhen this occurs the rear surface portion of the aperture formed in the outer end of the male connection member will be maintained substantially in contact with the pin member at all times.\nThis contact between such pin member and the rear surface portion of such aperture in the male connection member is absolutely necessary in this coupling arrangement because most of the articulated connecting members are manufactured as cast components. Consequently, in order for the manufacturer to achieve any significant reduction in the cost of this coupling device, such cast articulated connecting members receive very little, if any, finish-type machining which could provide the necessary and desirable dimensional control. In other words, these cast connecting members will generally be assembled into such coupling device as cast. Therefore, as a result of this cost-saving practice, it is quite often very difficult to provide an articulated coupling device which will be self-adjusting under the various wear conditions that will be encountered by such coupling device during in-track operation. Nevertheless, it is of the utmost importance to minimize the slack encountered in the various coupling connections during such in-track service for the reasons discussed above.\nOther prior art-type articulated coupling arrangements are taught in U.S. Pat. No. 3,716,146 and Canadian Patent number 1,231,078. However, the shortcomings which are overcome by the invention to be taught hereinafter are also found in these prior art references."}
-{"text": "An assembly as well as systems and methods for performing parallel crystallisation experiments are known from WO 02/06802. The known assembly comprises a microplate having multiple microwells having an opening at the top. Sealing of the wells is effected by O-ring seals around the top of each well, which are interpositioned between the microplate and a cover plate.\nThe known assembly is not sufficiently practical when conducting parallel experiments in high volumes, known as high throughput experimentation. In particular the known assembly is unsatisfactory when conducting parallel crystallisation experiments."}
-{"text": "The present invention relates to a differential amplifier having high gain at a low frequency signal input."}
-{"text": "\u201cHealth care professionals\u201d is used broadly here to refer to anyone who participates in the diagnosis or treatment of medical problems. For example, medical doctors, dentists, nurses, nurse-practitioners, medical technologists, physical therapists, and other health workers that assist in examination of patients, diagnosis, or treatment are all included by this term.\nA health care professional diagnoses an illness by collecting and evaluating information about the patient, then determining what disease or condition best fits the information. The information gathered from the patient usually is processed to reach a diagnosis by using a protocol learned during the professional's professional training and modified and updated by his or her medical experience. The protocol is an ordered process by which a health care professional ascertains information that allows the professional to rule out possible diseases until enough information is gathered to eliminate all but the diagnosed condition. Alternatively, the protocol may end when an appropriate treatment is identified. Recently, medical associations, health maintenance organizations, and hospitals, among others, have prescribed protocols. Employed health care professionals in particular are often subject to mandated protocols.\nOne problem in the field of medicine is how to improve diagnostic protocols to take into account advances in medical knowledge. A related problem is how to ensure that health care professionals update their skills to take advantage of advances in medical knowledge. Still another problem is how to expedite the diagnosis and treatment of certain conditions that should be treated quickly, so treatment can begin soon enough to be most effective.\nU.S. Pat. No. 6,095,973 discloses a data processing system and method for evaluating the treatment of chest pain patients in a medical facility.\nU.S. Pat. No. 6,029,138 discloses a decision support system for the selection of a diagnostic test or therapeutic intervention, which are both called \u201cstudies\u201d in that patent.\nU.S. Pat. No. 4,857,713 discloses a program for reducing hospital errors in the delivery of medications, goods, services or procedures in patient treatment.\nU.S. Pat. No. 5,732,397 describes an automated system for use in decision-making processes, which is said to improve the quality and consistency of decisions made.\nU.S. Pat. No. 5,772,585 discloses a common user interface to allow different medical personnel access to centralized files regarding patients.\nU.S. Pat. No. 5,832,450 describes an electronic medical record system that stores data about individual patient encounters in a convenient form.\nU.S. Pat. No. 5,845,255 describes an electronic prescription creation system for physician use that includes an adverse indication review and online access to comprehensive drug information including scientific literature.\nU.S. Pat. No. 5,911,132 discloses diagnosing and treating patient diseases using a epidemiological database containing medical, personal or epidemiological data relevant to a presented set of symptoms, test results, a diagnosis, etc.\nU.S. Pat. No. 5,915,240 discloses a context-sensitive medical lookup reference computer system for accessing medical information over a network.\nU.S. Pat. No. 5,924,074 discloses a medical records system that is said to create and maintain all patient data electronically.\nU.S. Pat. No. 5,953,704 discloses a system in which a user inputs information related to the health condition of an individual.\nU.S. Pat. No. 6,022,315 discloses a system and method for providing computerized, knowledge-based medical diagnostic and treatment advice to the general public over a telephone network or a computer network.\nThere is currently a need in the medical field for a system that communicates to a health care professional carrying out a diagnosis that a certain symptom, combination of symptoms, or other patient information recorded by the physician is associated with an increased risk of a missed medical care opportunity leading to a less favorable patient outcome. (A \u201cmedical care opportunity\u201d is defined as an opportunity to correctly or more quickly diagnose or treat the patient's condition and thus provide a better patient outcome.) Further, there is also a need in the medical field for a system for communicating to the health care professional special steps to take to avoid the missed medical care opportunity."}
-{"text": "The present invention relates generally to multiplexing of communications signals and more particularly to the wavelength division multiplexing of optical signals in fiber optics communications.\nWavelength Division Multiplexing (WDM) is a way of increasing the capacity of an optical fiber by simultaneously operating at more than one wavelength within a single optical fiber. Multiple optical signals of different wavelength are transmitted in the same direction over one strand of fiber, and the signals are later separated by wavelength at the distant end. In order to establish some standards for WDM operations, the International Telecommunication Union (ITU) has proposed successive systems of standardized frequencies to be used as channels for optical telecommunications, with each system incorporating more and more channels, usually with smaller and smaller wavelength separation between the channels. This system of channels is spoken of as the xe2x80x9cITU gridxe2x80x9d and currently includes 80 channels utilizing a wavelength range centered around 1550 nm (193,300 GHZ ) with a channel spacing of approximately 0.4 nm (50 GHZ). An ITU grid with 50 GHZ frequency spacing is shown in FIG. 1 with the fringe order pattern from a Fabry-Perot interferometer using an etalon of appropriate parameters superimposed on the grid. There has also been a system proposed which uses channel spacing of 25 GHZ (0.2 nm). It will be easily apparent that the smaller the spacing is between channels, the more sensitive a multiplexed channel system will be to fluctuations that will cause the channel frequencies to drift away from the target grid frequencies. In a system where the frequency spacing is 100 GHz, the variation of from the grid frequency of 0.01 nm (1.2 GHz) would be undesirable, but perhaps may not be disastrous. In a 50 GHz system, a 0.01 nm variation would impair the performance of the system, and in a 25 GHz system, such a deviation would likely be disabling.\nAn interference filter is a type of tool that is often used to separate multiple wavelengths of light that are included in a beam of light. A Fabry-Perot Interferometer is one type of interference filter that is often used for wavelength filtering and separation. Interference filters operate by providing a pair of mirrored surfaces with a cavity formed between them. Incident light undergoes multiple reflections between the mirrored surfaces, which typically reflect greater than 95% of the light at each surface. The incident and reflected wave interfere with each other constructively or destructively depending on their phase relationship. Where there is no phase difference between successive wave, constructive interference is produced and a maximum is produced in the transmission portion. Where the phase are 180 degrees out of phase, destructive interference occurs and a minimum is transmitted. A maximum occurs when the round trip optical path is some integer multiple of whole wavelengths, and also depends on the thickness of the cavity (d), the index of refraction of the cavity material (n), and the angle of incidence (0), which are related by the formula:\n2d n cosxcex8=mxcex,\nwhere m is an integer, often termed the order number and xcex is the wavelength of the light. The parallel rays of each wavelength are often focused by a lens in order to produce a familiar ring pattern. The result is a series of transmission peaks of separated wavelength. The separation distance between adjacent peaks is equally spaced when plotted with respect to inverse wavelength, and is called the Free Spectral Range (FSR).\nEtalons are special Fabry-Perot interferometers which have fixed spacing between the reflective surfaces, thus the thickness of the cavity d is therefore not subject to direct parallel variation. However, the etalon may be tilted, changing the angle of the etalon relative to the angle of incidence of the light beam, which thus increases the optical path length. This allows the etalon to be xe2x80x9ctunedxe2x80x9d over a limited range to alter the peak transmission wavelengths. FIGS. 2 and 3 show etalons as used in the prior art. FIG. 2 shows the incident light striking the etalon at an angle xcex8, and FIG. 3 shows the same etalon tuned at the same angle xcex8, thereby increasing the optical path length n to angle tune the etalon.\nHowever, angle tuning of an etalon introduces other problems caused by the insertion loss due to the variation in angle. Besides the difficulties of producing very tiny variations in angle, when the etalon is tuned at a small angle, the output beam can become oblong in shape, with non-uniform beam intensity distribution. As this angle increases, this effect becomes more pronounced. When optics are used to collect the output light from the etalon, a large insertion loss variation is often seen. This variation is typically from 1-4 dB.\nThe variation of etalon insertion loss also commonly causes the operating point of the output spectrum to shift by as much as 10 pm (0.01 nm=1.2 GHz). As discussed above, errors of this magnitude can seriously interfere with operation of systems which use 25 Ghz frequency separations and even with 50 GHz systems.\nThus, there is a great need for a etalon which is usable in a multi-channel wavelength locking system which does not produce such large variations in insertion loss, beam quality, and wavelength shift.\nAccordingly, it is an object of the present invention to provide a multi-channel wavelength locking system with better ITU setting repeatability.\nAnother object of the invention is to provide a multi-channel wavelength locking system with improved temperature stability.\nAnd another object of the invention is to provide a multi-channel wavelength locking system which produces smaller variations in insertion loss.\nA further object of the present invention is to provide a multi-channel wavelength locking system which produces smaller variations in beam quality.\nAn additional object of the present invention is to provide a multi-channel wavelength locking system which produces less wavelength shift.\nYet another object of the present invention is to provide a multi-channel wavelength locking system which can be used with wavelengths which are separated by as little as 25 GHz or less.\nA yet further object of the present invention is to provide a tunable etalon which has reduced manufacturing costs due to more relaxed dimensional tolerances in the parts.\nBriefly, one preferred embodiment of the present invention is a wavelength locker for providing wavelength-locked multi-channels signals, including a number of radiation sources providing beams of radiation at wavelengths separated by a predetermined spacing, each spaced wavelength being a channel which is connected to a multiplexer. One or more wavelength lockers produce equally-spaced spectral lines which are tunable wavelength locker includes one or more etalons which includes a gas-tunable medium having a variable optical index of refraction. The etalons may either be transmissive, or reflective with respect to which surface the major portion of the light exits. The etalons are tuned by varying the pressure or composition of the gas-tunable medium until the spectral lines of the etalon align with the ITU grid. The gas properties are then fixed by at least temporarily sealing the etalon enclosure.\nIn use in a multi-channel system, beam splitters divert a portion of the beams of radiation into the wavelength lockers. Detectors then receive the spectral line output of the wavelength lockers and detects shifts in wavelength of the spectral line output to generates control signals. The wavelength of each of the radiation sources is then adjusted in response to the control signals.\nThe output of the detector may also be normalized in order to isolate the variations which are attributable to shift in wavelength. The channels may be sequentially sampled and the outputs adjusted by a single wavelength locker or there may be separate wavelength lockers for each input channel.\nA method for locking the wavelengths of a plurality of radiation sources to a wavelength reference is also included.\nAn advantage of the present invention is that it provides improved ITU setting repeatability.\nAnother advantage of the present invention is that it provides improved temperature stability.\nAnd another advantage of the present invention is that it produces very small variations in insertion loss and beam quality.\nA further advantage of the present invention is that it can be used with wavelengths which are separated by as little as 25 GHz or less.\nA yet further advantage is that the present invention can provide a tunable etalon which has reduced manufacturing costs due to more relaxed dimensional tolerances in the parts.\nThese and other objects and advantages of the present invention will become clear to those skilled in the art in view of the description of the best presently known mode of carrying out the invention and the industrial applicability of the preferred embodiment as described herein and as illustrated in the several figures of the drawings."}
-{"text": "1. Field of the Invention\nThe present invention relates to a digital broadcast system for transmitting/receiving digital broadcast data, and a data processing method for use in the same.\n2. Discussion of the Related Art\nA variety of digital broadcast standards have been widely used throughout the world. For example, a vestigial sideband (VSB) transmission scheme has been used as a digital broadcast standard for the North America and the Republic of Korea. This VSB scheme is based on a single-carrier scheme, such that a reception (Rx) performance of a reception (Rx) system may deteriorate under poor channel environments. Specifically, a portable or mobile broadcast receiver requires much higher resistance to a channel variation or noise. So, if mobile service data is transmitted according to the VSB transmission scheme, a reception (Rx) performance becomes deteriorated."}
-{"text": "Other than referring to a store directory (if one exists) or finding a store sales clerk or a knowledgeable shopper, there is currently nothing available to help shoppers locate items in a store. Thus, many shoppers may search all over a store to locate one or more items, or may wait in a long line at a service counter to ask where particular items are located. This is a particularly difficult problem for handicapped shoppers who use motorized chairs. Such a shopper must press a button, or the like, in order to control movement of the chair. While searching for the location of store items, these shoppers may experience discomfort or fatigue from controlling the motorized chair. Further, shoppers who are hearing impaired may not be able to easily communicate with store employees and may not be able to effectively ask where an item is located. A quick and easy way to locate items in stores is needed to make shopping trips shorter and more productive."}
-{"text": "1. Field of the Invention\nThe present invention pertains to a pressure plate assembly for a friction clutch, which assembly includes a housing; a pressure plate connected to the housing for rotation in common around an axis of rotation which is free to move axially with respect to the housing; a stored-energy device which pretensions the pressure plate in the axial direction with respect to the housing; a wear-compensating device which acts in the path of support between the stored-energy device and at least one of the housing and the pressure plate and comprises at least one wear-compensating element, which is free to move when wear occurs and/or when wear is compensated; an indicator arrangement comprising at least one indicator opening in the housing and at least one indicator element on a wear-compensating element, where information concerning the amount of wear which has occurred can be derived from the position of the indicator element relative to the minimum of one indicator opening.\n2. Description of the Related Art\nU.S. Pat. No. 5,531,308 discloses a friction clutch in which a wear-compensating device can be used to ensure that the wear which occurs during the operation of the clutch, such as the wear of the friction linings of a clutch disk, for example, is compensated. For this purpose, two adjusting rings are provided, one of which is held nonrotatably with respect to a housing, whereas the other adjusting ring is free to rotate to compensate for wear. The two rings have corresponding ramp surfaces, so that, upon rotation of the second adjusting ring, the total axial length of the wear-compensating device consisting essentially of the two adjusting rings changes. On the adjusting ring which can rotate relative to the housing, an axial projection is provided, which engages in a curved, slot-like opening in the housing, i.e., in a bottom area of the housing. When the adjusting ring turns, this projection moves in the opening through which it passes. From the position which this projection occupies in the opening, it can be determined how far the rotatable adjusting ring has rotated up to that point, which is an indication in turn of the extent to which wear compensation has occurred and thus also of the amount of wear which has occurred.\nIt is the task of the present invention to improve a pressure plate assembly of the general type in question in such a way that, by the use of a simple design, the amount of wear which has occurred can be detected more easily.\nAccording to the invention, the indicator element at least partially fills at least one indicator opening in at least one state of wear.\nWith the design of a pressure plate assembly according to the invention, the geometry of the indicator element can be made essentially independent of the geometry of the indicator opening. This facilitates the design and also makes it possible for both the indicator element and the indicator opening to be designed with respect to their shape and/or configuration in such a way that, through the cooperation of these two assemblies, a very precise indication of the actual amount of wear which has occurred or of the degree of wear compensation which has occurred can be obtained.\nIt is possible, for example, for the width of the indicator element, i.e., the dimension essentially transverse to the direction of its movement with respect to the minimum of one indicator opening, to change over the length of the indicator element in the direction of movement. It is thus achieved that the width of the indicator element visible in the minimum of one indicator opening changes as the element moves past, so that, as the indicator element moves, the area of the indicator opening which it fills also changes, from which information can be derived concerning different states of wear.\nFor this purpose, it is possible, for example, for the width of the indicator element to change continuously, at least in certain of its areas. Alternatively or in addition, it is possible for the width of the indicator element to change in step-like increments. Although a step-like change in width produces a very clear-cut discontinuity in the indicator characteristic, an indicator element with a continuous change in width provides information of higher resolution concerning the state of wear.\nIt is also possible to provide a plurality of indicator openings in succession in the direction of the movement of the indicator element. In another, preferred embodiment of the pressure plate assembly according to the invention, it is possible, in conjunction with the provision of a plurality of indicator openings, for the indicator element to have an indicator section which fills essentially all of the indicator openings in one wear state and at least partially fills one indicator opening or essentially no indicator opening in at least one other wear state. It can therefore be provided, for example, that the indicator fills all of the indicator openings when the clutch is new and has not yet been affected by wear, and that, as the amount of wear increases, the indicator element moves successively along the indicator openings and thus opens various indicator openings one after the other.\nAs an alternative, it is also possible for the indicator element to have an indicator section which does not fill all of the indicator openings in any wear state. It can be provided in this case, for example, that, as a function of the amount of wear which has occurred, the indicator section completely fills one of the indicator openings or partially fills two of the indicator openings.\nIn a design of the pressure plate assembly which is very simple to realize, it is possible for the wear-compensating device to comprise an adjusting ring, which can rotate around the axis of rotation for the purpose of wear compensation, and for the indicator element to be free to move along with the adjusting ring constituting the wear-compensating element.\nAn alternative embodiment of the pressure plate assembly according to the invention can be designed in such a way that the wear-compensating device comprises: an arresting element which, upon the occurrence of wear, is able to move with respect to the pressure plate assembly in correspondence with the amount of wear which has occurred; an essentially wedge-shaped slider element, which is installed between the pressure plate assembly and the arresting element and which, upon movement of the arresting element with respect to the pressure plate assembly, is free to move under the pretensioning effect in correspondence with the extent of the movement of the arresting element with respect to the pressure plate assembly in order to keep the arresting element in the position relative to the pressure plate assembly to which it has moved as a result of wear; and an adjusting ring, which is able to rotate around the axis of rotation to perform a wear-compensating movement, where the adjusting rotation of the adjusting ring is limited by the slider element and/or by the arresting element, where the indicator element is able to move along with the slider element forming the wear-compensating element.\nIn addition, the pressure plate assembly according to the invention can be designed so that the stored-energy device is located in an area between the indicator element and the minimum of one indicator opening and has a see-through opening associated with the minimum of one indicator opening. To avoid the introduction of additional openings in the stored-energy device, it is possible, for example, for the see-through opening to be formed by an area between two of the spring tongues of the stored-energy device.\nIt is also possible, according to the invention, for the indicator element to be connected to the wear-compensating element by connecting sections provided on the indicator element. It is also possible for the indicator element to have a contact area, by means of which it is supported on a surface of the wear-compensating element facing the axis of rotation. It is possible in this way to prevent the indicator element from shifting position as a result of centrifugal force during rotational operation.\nThe present invention also pertains to a friction clutch which contains a pressure plate assembly according to the invention.\nOther objects and features of the present invention will become apparent from the following detailed description considered in conjunction with the accompanying drawings. It is to be understood, however, that the drawings are designed solely for purposes of illustration and not as a definition of the limits of the invention, for which reference should be made to the appended claims. It should be further understood that the drawings are not necessarily drawn to scale and that, unless otherwise indicated, they are merely intended to conceptually illustrate the structures and procedures described herein."}
-{"text": "The structure of the growth hormone release inhibiting factor, somatostatin, has been determined by Brazeau et al., Science, 179, 77(1973). Several techniques for synthesizing somatostatin have been reported in the literature, including the solid phase method of Rivier, J.A.C.S. 96, 2986(1974) and the solution methods of Sarantakis et al., Biochemical Biophysical Research Communications 54, 234(1973) and Immer et al., Helv. Chim. Acta, 57, 730(1974)."}
-{"text": "Dinettes are common features in recreational vehicles. A typical dinette includes a pair of bench seats and a table supported by one or more table posts or legs. Because space is limited in recreational vehicles, dinettes often can be converted into a bed by lifting the tabletop from its support posts, removing the support posts, and inserting the tabletop between the dinette seats on small ledges on the seats. The problem with conventional convertible dinettes is that the support posts or legs must be stowed when used as a bed."}
-{"text": "1. Field of the Invention\nThe present invention relates to vapor generators and, more particularly, to a vapor generator fluid flow control system, permitting controlled fuel, air and water flow and pressure regulation.\n2. History of the Prior Art\nVapor generators of the kind in which a fuel-air mixture is combusted in the direct presence of feed water to produce a useful mixture of steam and non-condensibles are known. Such vapor generators are shown in U.S. Pat. No. 4,211,071 issued to the assignee of the present invention. In accordance with that invention, a vapor generator is provided in which several interrelated means are employed to improve the quality of combustion in the generator so that a product steam substantially free of carbon monoxide results. In effecting this result, means are provided for dividing the air feed into two parts which are both delivered to the generator for combustion. Such methods and apparatus are effective in producing a well-mixed stoichiometric mixture of fuel and air to provide completeness of combustion and reducing production of carbon monoxide to extremely low levels. Such an accomplishment is a marked advance over the prior art and illustrates the emphasis placed on controlled flow rates and mixture conditions.\nThe prior art is also replete with apparatus manifesting high temperature product vapor having high carbon monoxide content. Such conditions are objectionable for many applications and dangerous for some of them. High carbon monoxide production is traceable to incomplete combustion. This is, in turn, traceable in part to difficulties in maintaining a stable lean flame and, in part, to excessive quenching of the flame through direct radiation and convective contact between the flame and the feed water. Once a stoichiometric mixture is obtained, it is, therefore, necessary to maintain the combustion level without varying any of the three major constituents of fuel, air, or water. Variations would affect the combustion and the carbon monoxide production therefrom.\nOther problems occur in certain ones of the prior art vapor generators when they are operated at low pressures. Conventionally, as shown in U.S. Pat. No. 4,288,978, a vapor generator may be especially adapted for low pressure operation by the provision of certain features such as a pilot burner for striking a stable flame. It may be seen that low pressure operation, likewise, requires a specific reduction in the air, fuel and water input for a low volume output from the vapor generator. Disproportionate variances in any of the three constituents will seriously affect the combustion level and the resulting heat and vapor content of the downstream product.\nFew prior art systems effectively address dual output vapor generator operation. The ability to fine tune the fuel, air and water flow for a single vapor generator output is critical enough. Once this critical equilibrium is reached, it is tantamount to starting over to vary the generator operation to another flow. Varying the generator volume necessitates adjustment in the air, fuel and water flow rates which will directly affect both the heat and the water vapor level of the combustion products. When such generators are used in industrial applications necessitating precise control of heat and vapor, these variations are not tolerable. One such application is the curing of concrete where the temperature and moisture level of the vapor products have a direct effect on the curing process and resulting structural condition of the cast product.\nThe use of steam for concrete curing is not novel in and of itself. Steam has been produced from conventional boilers for such methods and processes, but the cost of a boiler operation is much higher than that of a vapor generator type device. Moreover, the temperature and vapor level of the steam product can be more closely controlled than with a conventional boiler. Moreover, super heated vapor can be provided for specific types of aggregate to be cured. The flexibility of the vapor generator also permits substantially instantaneous operation as compared to the threshold period for a boiler method. The obvious disadvantage is the substantially single flow rate afforded by most conventional vapor generators. Since any number of kilns may be run at one time for particular production schedules, excess vapor from a single flow rate vapor generator system would, thereby, be wasted.\nConventional techniques for varying the output of standard vapor generators includes reducing the rate of combustion and water flow while maintaining a constant air flow rate. The advantage of such a system is simplicity and cost in that reliable variations in air flow rates have been, to date, difficult to attain. The paramount disadvantage to the continuous air flow volume is the excess air ingressing into the kiln and the adverse effects to the curing process which affects the ultimate humidity level within the kiln and imparts non-uniform cooling characteristics. Variations in the air flow volume have, to date, been addressed by multiple speed blowers and throttle systems which either decrease the air intake to the blower or exhaust therefrom. Unfortunately, the relatively large motors and blower units necessary for vapor generator volumes generate large amounts of heat which the air volume dissipates. When the air volume is throttled for constant motor speed, heat dissipation is inhibited and variations in flow rate result as well as having degenerative effects on the motor and the blower unit. The obvious alternative to such flow problems is a multiple speed motor, but the cost and availability for such systems are generally disproportionate to the vapor generator construction.\nDual output vapor generator applications also require dual flow feedwater systems which may be accurately set in preselected flow configurations as well as precisely maintained. Precise fluid flow is an integral element of the aforesaid combustion efficiency particularly in stoichiometric mixtures. The primary area of fluid flow consideration in such vapor generator systems is in the water flow controlling network. Prior art flow control devices though capable of preselect flow volumes are generally not sensitive to variations in flow pressure which affects the flow volume.\nIt would be an advantage therefore to overcome the problems of the prior art by providing a variable output vapor generator having a constant volume fluid flow system which affords automatic operation between high and low fluid flow rates and which is sensitive to fluctuations in fluid flow pressure. One approach to dual output vapor generators and a discussion of prior art problems associated therewith is discussed in co-pending U.S. patent application Ser. No. 554,780 assigned to the assignee of the present invention. The method and apparatus of the present invention comprises an advanced system overcoming the problems set forth above by providing a combination of a pressure responsive valve and fine adjustment valve in parallel flow communication. This parallel flow system is itself aligned in parallel flow communication with a second matching system adapted for selective actuation for high and low volume operation. In this manner both the high volume and low volume flow rates within the water flow system of a vapor generator may be pre-adjusted for automatic actuation as needed by demand conditions. This has been done in a configuration facilitating automatic temperature control and which is pressure sensitive for accurate volumetric control without adversely affecting the efficiency of the generator system."}
-{"text": "The present invention relates generally to medical devices and more particularly to delivery systems for self-expandable medical devices.\nSelf-expanding medical devices are used by physicians to treat numerous conditions using minimally invasive procedures. Examples of self-expanding medical devices include stents, stent-grafts, filters, valves, etc. Typically, self-expanding medical devices are made from an elastic structure that may be compressed into a low profile state that can be passed through vessels in a patient with minimal trauma. Once at the desired treatment site, the self-expanding medical device is released and self-expands like a spring until it contacts a tissue wall which prevents further expansion. Common materials that are used in self-expanding medical devices include nitinol and stainless steel, although other materials are also possible.\nOne type of self-expanding medical device that has become especially common is intraluminal stents. Stents are used to treat various organs, such as the vascular system, colon, biliary tract, urinary tract, esophagus, trachea and the like. For example, stents are commonly used to treat blockages, occlusions, narrowing ailments and other similar problems that restrict flow through a passageway. One area where stents are commonly used for treatment involves implanting an endovascular stent into the vascular system in order to improve or maintain blood flow through narrowed arteries. However, stents are also used in other treatments as well, such as the treatment of aneurysms. Stents have been shown to be useful in treating various vessels throughout the vascular system, including both coronary vessels and peripheral vessels (e.g., carotid, brachial, renal, iliac and femoral). In addition, stents have been used in other body vessels as well, such as the digestive tract.\nThe use of stents in coronary and peripheral vessels has drawn particular attention from the medical community because of the growing number of people each year that suffer from vasculature problems associated with stenosis (i.e., narrowing of a vessel). This has led to an increased demand for medical procedures to treat such problems. The widespread frequency of heart problems and other vasculature problems may be due to a number of societal changes, including the tendency of people to exercise less and the prevalence of unhealthy diets, in conjunction with the fact that people generally have longer life spans now than previous generations. Stents have become a popular alternative for treating vascular stenosis because stenting procedures are considerably less invasive than conventional procedures. For example, stenosis of the coronary arteries was traditionally treated with bypass surgery. In general, bypass surgery involves splitting the chest bone to open the chest cavity and grafting a replacement vessel onto the heart to bypass the blocked, or stenosed, artery. However, coronary bypass surgery is a very invasive procedure that is risky and requires a long recovery time for the patient. Vascular stents are also being more widely used to treat many different peripheral arteries due to the minimally invasive nature of stenting procedures. To address the growing demand for minimally invasive medical procedures for the treatment of coronary arteries, peripheral arteries and other passageway problems, the medical community has begun to turn away from conventional invasive procedures like bypass surgery and increasingly the treatment of choice now involves a variety of stenting procedures.\nMany different types of stents and stenting procedures are possible. In general, however, stents are typically designed as tubular support structures that may be inserted percutaneously and transluminally through a body passageway. Traditionally, stents are made from a metal or other synthetic material with a series of radial openings extending through the support structure of the stent to facilitate compression and expansion of the stent. Although stents may be made from many types of materials, including non-metallic materials, common examples of metallic materials that may be used to make stents include stainless steel, nitinol, cobalt-chrome alloys, amorphous metals, tantalum, platinum, gold and titanium. Typically, stents are implanted within a passageway by positioning the stent within the area to be treated and then expanding the stent from a compressed diameter to an expanded diameter. The ability of the stent to expand from a compressed diameter makes it possible to thread the stent to the area to be treated through various narrow body passageways while the stent is in the compressed diameter. Once the stent has been positioned and expanded at the area to be treated, the tubular support structure of the stent contacts and radially supports the inner wall of the passageway. As a result, the implanted stent mechanically prevents the passageway from narrowing and keeps the passageway open to facilitate fluid flow through the passageway.\nStents can generally be characterized as either balloon-expandable or self-expanding. Traditionally, balloon-expandable stents have been used most often in coronary vessels than in peripheral vessels because of the deformable nature of these stents. One reason for this is that peripheral vessels tend to experience frequent traumas from external sources (e.g., impacts to a person's arms, legs, etc.) which are transmitted through the body's tissues to the vessel. In the case of peripheral vessels, there is an increased risk that an external trauma could cause a balloon-expandable stent to plastically deform in unexpected ways with potentially severe and/or catastrophic results. However, in the case of coronary vessels, this risk is minimal since coronary vessels rarely experience traumas transmitted from external sources.\nSelf-expanding stents are increasingly used and accepted by physicians for treating a variety of ailments. Self-expanding stents are usually made of shape memory materials or other elastic materials that act like a spring. Typical metals used in this type of stent include nitinol and 304 stainless steel. A common procedure for implanting a self-expanding stent involves a two-step process. First, the narrowed vessel portion to be treated is dilated with a balloon but without a stent mounted on the balloon. Second, a stent is implanted into the dilated vessel portion. To facilitate stent implantation, the stent is installed on the end of an inner catheter in a compressed, small diameter state and is usually retained in the small diameter by inserting the stent into an outer sheath at the end of the catheter. The stent is then guided to the balloon-dilated portion and is released from the inner catheter by pulling the outer sheath away from the stent. Once released from the outer sheath, the stent radially springs outward to an expanded diameter until the stent contacts and presses against the vessel wall. Traditionally, self-expanding stents have been more commonly used in peripheral vessels than in coronary vessels due to the shape memory characteristic of the metals that are used in these stents. One advantage of self-expanding stents for peripheral vessels is that traumas from external sources do not permanently deform the stent. Instead, the stent may temporarily deform during an unusually harsh trauma but will spring back to its expanded state once the trauma is relieved. Self-expanding stents, however, are often considered to be less preferred for coronary vessels as compared to balloon-expandable stents. One reason for this is that balloon-expandable stents can be precisely sized to a particular vessel diameter and shape since the ductile metal that is used can be plastically deformed to a desired size and shape. In contrast, self-expanding stents are designed with a particular expansible range. Thus, after being implanted, self-expanding stents continue to exert pressure against the vessel wall.\nCommonly, delivery systems for self-expanding medical devices have a handle arrangement that remains outside of the patient's body during the deployment procedure. One portion of the handle is typically connected to an inner catheter upon which the self-expanding medical device is mounted, and another portion of the handle is typically connected to an outer sheath which restrains the self-expanding medical device in the compressed state. When the distal end of the delivery system is positioned within the patient's body at the intended treatment site, the physician actuates the handle by moving the two portions relative to each other so that the outer sheath is withdrawn from the self-expanding medical device and inner catheter. As a result, the self-expanding medical device expands outward away from the inner catheter. The handle may then be pulled by the physician to withdraw the inner catheter and outer sheath from the patient's body, while leaving the self-expanding medical device implanted in the body.\nPrecise placement of self-expanding medical devices is a concern in most medical procedures. However, precise placement can be more difficult with certain delivery systems due to their design, shape and other factors. Precise placement of self-expanding medical devices is generally a function of the relative movement and placement between the delivery system handle and the patient's body, and the relative movement between the portions of the handle connected to the inner catheter and outer sheath during deployment. A lack of control over any part of this system can result in inaccurate placement of a self-expanding medical device, and thus, less than desirable treatment of the medical condition being treated.\nAccordingly, the inventor believes it would be desirable to provide a new delivery system for self-expanding medical devices."}
-{"text": "1. Field of the Invention\nThe present invention relates to an electronic component feeding apparatus for conveying electronic components by intermittent movement of a belt.\n2. Description of the Prior Art\nJapanese Patent Application Laid-Open No. Heisei6-232596 has disclosed an apparatus for conveying chip components using a belt.\nThis apparatus has an endless belt wound around a pair of pulleys, a ratchet mechanism for intermittently moving a belt through a predetermined distance by driving one of the pulleys, and a lever mechanism for transmitting the power to the ratchet mechanism. When the power is transmitted to the ratchet mechanism by moving an operating portion of the lever mechanism in a predetermined direction, one pulley is turned through a predetermined angle in a predetermined direction by the ratchet mechanism, and the belt is moved through a predetermined distance in a predetermined direction by the turning of the pulley. Chip components, arranged in an orderly way on the belt, are conveyed in the same direction by the movement of the belt, and the foremost chip component is fed to a predetermined takeout position. The foremost chip component fed to the takeout position is taken out by using a suction nozzle. When the operating portion of the lever mechanism is moved again, the succeeding chip component is fed to the takeout position in the same way as described above.\nThe apparatus of this type has widely been used for a component mounter, component mounting line, and the like. In recent years, there has been demanded a mechanism which can respond to high-speed component takeout cycle, specifically a mechanism which can take out components in a very short cycle as short as 0.1 sec.\nHowever, on the aforesaid apparatus, since the operating portion of the lever mechanism is directly connected to an driving portion of the ratchet mechanism, the speed at which the belt moves is substantially proportional to the speed at which the operating portion of the lever mechanism is moved. That is to say, if the operation speed of the operating portion of the lever mechanism is increased to respond to the high-speed component takeout cycle, the belt movement speed increases substantially proportionally, so that the chip components on the belt cannot follow the belt initial speed. Thereby, a slip occurs between the belt and the chip components on the belt, so that only the belt moves while the chip component is left, with the result that components cannot be conveyed satisfactorily by means of the belt."}
-{"text": "This invention relates generally to an implantable containment apparatus made of selectively permeable material. In particular, the implantable containment apparatus is used to contain a therapeutical device, such as a drug delivery device, a cell encapsulation device, or a gene therapy device. A therapeutical device can be easily placed and replaced in an apparatus of the present invention without damaging tissues associated with the selectively permeable material of the apparatus.\nVarious implantable therapeutical devices, such as drug delivery, gene therapy, and cell encapsulation devices, have been disclosed over the years. A common feature of most of these devices is the use of selectively permeable, or semi-permeable, membranes to construct all or part of the device. These membranes contain their respective therapeutic agents and delivery systems within the particular device while being permeable to the desired therapeutical product. For cell encapsulation devices, the membranes are also permeable to life sustaining substances and to cellular waste products.\nWhen implanted in a recipient, the typical biological response by the recipient to most of these therapeutical devices is the formation of a fibrotic capsule around the device. With most drug delivery and gene therapy devices, this can limit the performance of the device, particularly when the therapeutic agent has a short half-life. For cell encapsulation devices, a fibrotic capsule encasing the device most often deprives the encapsulated cells of life sustaining exchange of nutrients and waste products with tissues of a recipient. The result is usually fatal to the encapsulated cells. Furthermore, a fibrotic capsule encasing a therapeutical device usually makes surgical retrieval of the device difficult.\nWhen certain therapeutical devices are implanted in a recipient, predominantly vascular tissues of the recipient can be stimulated to grow into direct, or near direct, contact with the device. On one hand, this is desirable because the therapeutical product of the device can then be delivered directly to the circulation of the recipient through the vascular tissues that are in contact with the device. On the other hand, this is undesirable because once vascular tissues of a recipient have grown in contact with one of these implantable therapeutical devices, removal of the device requires surgical dissection of the tissues to expose and remove the device. Surgical dissection of vascular tissues, particularly capillary tissue, can often be a difficult and painful procedure. Whether encased in a fibrotic capsule or surrounded with vascular tissue, the problem of retrieving these implanted devices is a considerable drawback of the devices.\nFor cell encapsulation devices, an alternative to retrieving and replacing the entire device in a recipient is to retrieve and replace the cells contained in the device. U.S. Pat. No. 5,387,237, issued to Fournier et al., is a representative example of a cell encapsulation device that has at least one opening into the device through which cells can be introduced and removed. Cells are introduced and removed in this, and other similar devices, as a suspension or slurry. Since most cell encapsulation devices are intended to correct a metabolite deficiency in a recipient caused by dysfunction or failure of certain of the recipient\"\"s cells, tissues, or organs, the source of the replacement cells is rarely the recipient. In a situation where non-autologous cells are used in this type of cell encapsulation device, the problem of contaminating a recipient with the foreign cells during loading, removal, or refilling of the device is ever present. One solution to this contamination problem would be to enclose the cells in a container that can be placed, removed, and replaced in a device as a unit.\nA retrievable cell encapsulation envelope enclosed in an implantable permselective membrane for use as an artificial endocrine gland is disclosed in U.S. Pat. No. 4,378,016 issued to Loeb. The Loeb device comprises a housing made of an impermeable hollow stem and a permselective membrane sack. The hollow stem has a distal end defining an extracorporeal segment, a percutaneous segment in the mid-region, and a proximal end defining a subcutaneous segment. The sack is adapted to receive an envelope containing hormone-producing cells and has an access opening that is coupled to the proximal end of the hollow stem. In a preferred embodiment, the cell containing envelope is in the form of a flexible collar. The flexible collar is partially collapsible to allow for easier placement and replacement of the envelope in the sack. Once in place, the flexible collar also provides a snug fit between the envelope and the sack. Placement and replacement of a cell containing envelope in the sack portion is accomplished manually with forceps, or the like. Retrieval of the envelope from the sack can be aided with a guidewire attached to the envelope. In one embodiment of the Loeb device, the sack has openings at both ends that are implanted percutaneously. In this embodiment, the cell containing envelope may be inserted or removed through either end of the device.\nThe housing of the Loeb device is surgically implanted in a recipient through the abdominal wall so the distal end of the stem protrudes from the recipient, the proximal end of the stem resides subcutaneously with respect to the abdominal wall, and the sack portion is placed in the peritoneal cavity surrounded in peritoneal fluid. According to Loeb, the sack allows hormones, nutrients, oxygen, and waste products to flow in and out of the sack while preventing bacteria from entering the patient. The sack and the envelope are said by Loeb to be permeable to nutrients and hormones, but impermeable to the hormone-producing cells and immune response bodies. Upon implantation of the device in a patient, the cells contained therein are said to take over the function of the corresponding natural gland, sense the amount of hormone needed, and produce the correct amount of the desired hormone.\nImplanted cell encapsulation devices, particularly those intended as an artificial endocrine gland, usually require a high rate of flux of nutrients and waste products between the encapsulated cells in the device and tissues of the recipient. Having a cell encapsulation device in close, or direct, association with a vascular structure usually provides the highest rate of nutrient and waste product flux for such a device. Loeb does not teach the value of vascularization of the sack portion of the housing, however. Nor is the Loeb device implanted in a part of the body that is particularly vascularized.\nBrauker et al. disclose a cell encapsulation device in U.S. Pat. No. 5,314,471 that requires close association of host vascular structures with the device. According to Brauker et al., xe2x80x9cconventional implant assemblies and methodologies usually fail to keep the implanted cells alive long enough to provide the intended therapeutic benefit.xe2x80x9d Cell death in these implanted devices is said by Brauker et al. to be due in large part to an ischemia imposed on the cells during the first two weeks following implantation. Brauker et al. conclude that xe2x80x9cthe cells die because conventional implant assemblies and methodologies themselves lack the innate capacity to support the implanted cells\"\" ongoing life processes during the critical ischemic period, when the host\"\"s vascular structures are not nearby.xe2x80x9d Brauker et al. state that in order for implanted cells to survive and function on a long term basis, the host must grow new vascular structures in association with the device. Brauker et al. note that a host will not naturally provide new vascular structures to an implanted cell encapsulation device. According to Brauker et al., the host must be stimulated by the implant assembly itself to grow new vascular structures close to the cell encapsulation device. Angiogenic stimulus can be provided by angiogenic factors applied to the cell boundary of the Brauker et al. device or by certain cell types encapsulated in the device. Growth of vascular tissue in association with a device of Brauker et al. will make removal of the device from a host difficult, however,\nAn implantable containment apparatus made of selectively permeable polymeric material that permits a therapeutical device, such as a drug delivery, a gene therapy, or a cell encapsulation device, to be placed and replaced in a recipient without damaging or disturbing tissues associated with the selectively permeable polymeric material would be useful. Such an apparatus that becomes closely associated with vascular structures without the need to supply angiogenic factors to induce the close vascularization would also be useful. A method of easily placing and replacing a therapeutical device in an implantable containment apparatus of the present invention would be further useful.\nThe present invention is directed to an implantable containment apparatus for a therapeutical device, such as a drug delivery device, a cell encapsulation device, or a gene therapy device. The apparatus is primarily made of selectively permeable material. The selectively permeable material permits flux, or exchange, of solutes between a therapeutical device contained in the apparatus and tissues of a recipient, while excluding cells from growing beyond a desired point through the material. When an apparatus is implanted in a recipient, various tissues from the recipient grow to associate with the apparatus. These tissues grow next to, or partially through, the exterior surface of the apparatus. It is preferred that vascular tissue is the predominant tissue that grows to associate with an apparatus of the present invention. Once growth of a recipient\"\"s tissues in association with an apparatus of the present invention has occurred, a therapeutical device can then be easily placed and replaced in the apparatus without damaging or disturbing the tissues associated with the selectively permeable material of the apparatus.\nThis is preferably accomplished by providing an implantable containment apparatus for a therapeutical device comprising a selectively permeable microporous polymeric material in the form of a tube wherein the tube comprises an exterior surface, an interior surface that defines a luminal space of substantially uniform diameter, and an access means at one end of the tube through which a therapeutical device is insertable into the luminal space of the tube wherein, once a therapeutical device is inserted into the luminal space of the tube, the therapeutical device is retained within the tube wherein biochemical and therapeutic substances having a molecular weight up to about 5,000,000 MW diffuse across the thickness of the tube between the contents of the therapeutical device contained therein and tissues of a recipient wherein the therapeutical device is removable from the tube through the access means of the tube and wherein the tube is refillable with a therapeutical device through the access means of the tube. The apparatus can also have an access means at each end of the tube. In this embodiment, a therapeutical device can be inserted and removed from the luminal space of the tube through either access means in the tube. Furthermore, in this embodiment, a fluid stream can be established through both access means of the tube that can then be used to flush a therapeutical device in and out of the luminal space of the tube.\nAccordingly, the present invention is also directed to a method in which a therapeutical device, such as a drug delivery device, a gene therapy device, or a cell encapsulation device, is easily inserted into, removed from, and replaced in a containment apparatus of the present invention as a unit with a fluid stream. The method involves repeatedly filling and emptying an implantable containment apparatus tube with a therapeutical device.\nOther features and advantages of the invention will become apparent upon review of the following specification, drawings, and claims."}
-{"text": "Visual inspection is widely used in the tire manufacturing process and even more commonly relies on the skill of the operators tasked with checking for the absence of visible imperfections at the surface of the tires in order to ensure the compliance thereof.\nHowever, with the advances in the computing power of computer-based means, tire manufacturers are developing automatic inspection means to assist the operators tasked with the visual inspection. To this end, it is possible to use an inspection device comprising lighting means and cameras which are positioned in such a way as to scan the exterior and interior zones of the lateral beads and of the tread of the tire that is to be inspected. The viewing field of each camera is angularly limited. In order to obtain complete images of the inside and of the outside of the tire, the tire has to be turned about its axis with respect to the lighting means and with respect to the cameras. The digital images obtained are then processed and compared against reference images in order to determine whether there might be any surface and appearance anomalies in the tire. For further details, reference may for example be made to patent applications EP-A2-1 959 227, EP-A1-2 023 078 and EP-A1-2 172 737.\nIn order to carry out such an inspection with a good degree of precision, use is generally made of an image acquisition device equipped with a laser lighting means that allows a line of light to be projected onto the tire and with a matrix camera capable of capturing the light reflected off the tire and which is oriented at a triangulation angle. The triangulation angle is the angle formed between the optical axis of the laser lighting means and the optical axis of the camera.\nIn order to be able to detect very small imperfections at the surface of the tire, the triangulation angle generally chosen is relatively large. As a result, the overall size of the image acquisition device is large.\nNow, in order to be able to capture a complete image of at least half the interior surface of the tire in a single rotation thereof, it is necessary to install a plurality of image acquisition devices penetrating the interior space of the tire.\nGiven the overall size of each image acquisition device it may therefore be particularly tricky if not to say impossible to achieve this installation for certain tire sizes, for example for tires in the current size range for passenger vehicles.\nThe present invention seeks to overcome this disadvantage."}
-{"text": "Lithographic printing typically involves the use of a so-called printing master such as a printing plate which is mounted on a cylinder of a rotary printing press. The master carries a lithographic image on its surface and a print is obtained by applying ink to said image and then transferring the ink from the master onto a receiver material, which is typically paper. In conventional lithographic printing, ink as well as an aqueous fountain solution (also called dampening liquid) are supplied to the lithographic image which consists of oleophilic (or hydrophobic, i.e. ink-accepting, water-repelling) areas as well as hydrophilic (or oleophobic, i.e. water-accepting, ink-repelling) areas. In so-called driographic printing, the lithographic image consists of ink-accepting and ink-abhesive (ink-repelling) areas and during driographic printing, only ink is supplied to the master.\nPrinting masters are generally obtained by the so-called computer-to-film method wherein various pre-press steps such as typeface selection, scanning, color separation, screening, trapping, layout and imposition are accomplished digitally and each color selection is transferred to graphic arts film using an image-setter. After processing, the film can be used as a mask for the exposure of an imaging material called plate precursor and after plate processing, a printing plate is obtained which can be used as a master.\nA typical photosensitive printing plate precursor for computer-to-film methods comprises a hydrophilic support and an image-recording layer which includes UV-sensitive compositions. Upon image-wise exposure of a negative-working plate, typically by means of a film mask in a UV contact frame, the exposed image areas become insoluble and the unexposed areas remain soluble in an aqueous alkaline developer. The plate is then processed with the developer to remove the diazonium salt or diazo resin in the unexposed areas. So the exposed areas define the image areas (printing areas) of the printing master, and such printing plate precursors are therefore called \u2018negative-working\u2019. Also positive-working materials, wherein the exposed areas define the non-printing areas, are known, e.g. plates having a novolac/naphtoquinone-diazide coating which dissolves in the developer only at exposed areas.\nIn addition to the above- photosensitive materials, also heat-sensitive printing plate precursors have become very popular. Such thermal materials offer the advantage of daylight-stability and are especially used in the so-called computer-to-plate method wherein the plate precursor is directly exposed, i.e. without the use of a film mask. The material is exposed to heat or to infrared light and the generated heat triggers a (physico-)chemical process, such as ablation, polymerization, insolubilization by cross-linking of a polymer or by particle coagulation of a thermoplastic polymer latex, and solubilization by the destruction of intermolecular interactions.\nThermal plates which require no processing are also known; such plates are typically of the so-called ablative type, i.e. the differentiation between hydrophilic and oleophilic areas is produced by heat-induced ablation of one or more layers of the coating, so that at exposed areas a surface is revealed which has a different affinity towards ink or fountain than the surface of the unexposed coating. A major problem associated with ablative plates, however, is the generation of ablation debris which may contaminate the electronics and optics of the exposure device and which needs to be removed from the plate by wiping it with a cleaning solvent, so that ablative plates are often not truly processless. Ablation debris which is deposited onto the plate's surface may also interfere during the printing process.\nOther thermal plates that require no processing are described in U.S. Pat. No. 5,855,173, U.S. Pat. Nos. 5,839,369 and 5,839,370 where a method relying on the image-wise hydrophilic-hydrophobic transition of a ceramic such as a zirconia ceramic and the subsequent reverse transition in an image erasure step. This image-wise transition is obtained by exposure to infrared laser irradiation at a wavelength of 1064 nm at high power (the average power is 1 W to 50 W and the peak power lies between 6 kW and 100 kW) which induces local ablation and formation of substoichiometric zirconia. U.S. Pat. No. 5,893,328, U.S. Pat. No. 5,836,248 and U.S. Pat. No. 5,836,249 disclose a printing material comprising a composite of zirconia alloy and \u03b1-alumina which can be imaged using similar exposure means to cause localized \u201cmelting\u201d of the alloy in the exposed areas and thereby creating hydrophobic/oleophilic surfaces. A similar printing material containing an alloy of zirconium oxide and Yttrium oxide is described in U.S. Pat. No. 5,870,956. The high laser power output required in these prior art methods implies the use of expensive exposure devices.\nAnother type of processless plates are printing plates based on a so-called \u201cswitching\u201d reaction where a hydrophilic surface is irreversibly changed into an oleophilic surface or vice versa by imagewise exposure. EP 652 483 for example, describes a positive working printing plate based on an acid catalyzed cleavage of acid-labile groups pendant from a polymer backbone. EP 200 488 and U.S. Pat. No. 4,081,572 describe negative working plates where a hydrophilic/hydrophobic conversion is obtained by a chemical reaction upon imagewise exposure to heat. Other examples of processless plates are based on the thermally induced rupture of microcapsules and the subsequent reaction of the microencapsulated oleophilic materials (isocyanates) with functional (hydroxyl-)groups on cross-linked hydrophilic binders (U.S. Pat. No. 5,569,573; EP 646 476; WO94/2395; WO98/29258).\nU.S. Pat. No. 6,582,882 describes an imaging element comprising a graft copolymer having a hydrophobic backbone and a plurality of pendant hydrophilic groups or a plurality of pendant groups comprising hydrophilic and hydrophobic segments. Upon exposure of the imaging element to thermal energy, the exposed areas become less soluble in a developer than the unexposed areas.\nU.S. Pat. No. 6,362,274 describes grafted copolymers comprising three sequences: one sequence for anchoring on solid particles such as pigments and fillers, one hydrophobic sequence and one hydrophilic sequence for using the copolymers in aqueous and/or organic medium. The disclosed copolymers are of particular interest in a wide range of paint formulations; there is no reference in the cited prior art document to lithographic printing plates.\nNone of the prior art discloses the heat-sensitive copolymer of the present invention in lithographic printing plates."}
-{"text": "The present invention relates to a device for auto-adaptative direction and polarization filtering of radio waves received on a network of aerials coupled to a receiver.\nThe invention is applied in particular to the construction of an interference elimination device for receivers of electromagnetic waves travelling through the HF ionospheric channel.\nIt is known that a tactical high frequency reception aerial must in particular enable any user to establish radio connections with various transmitting stations, and to simultaneously hear the transmissions from the other users of the high frequency channel.\nBearing in mind the congestion of the high frequency channel, the receiver must be able to reject the signals coming from sources of interference, be they intentional or not.\nHowever, the known tactical high frequency adaptative aerial networks, generally formed from several suitably spaced vertical whips, do not enable exploitation of the polarization of the ionospheric waves. Whilst the performance levels expected from an adaptative aerial formed from three orthogonal dipoles, which does not have the shortcomings of the whips, are described in an article by RTE COMPTON entitled \"The tripole antenna: an adaptative array with a full polarization flexibility\" (IEEE trans antenne and propagation. Vol.: AP 29 No. 6 Nov. 81), the aerial which is described therein has the tactical disadvantage of having to be placed at the top of a mast so that the radiation patterns of the horizontal dipoles are correct. Moreover, as it is technologically very difficult to suppress the high frequency currents circulating in the external parts of the coaxial cables connecting the dipoles of the aerial to the receiver, the performance levels expected of this aerial configuration are, in practice, not obtained."}
-{"text": "Currently, Chinese herbal injections for treatment of cardio-cerebral vascular diseases and fundus diseases include Ixeris Sonchifolia Hance injection, which for clinical application is a brown-yellow transparent liquid derived from Ixeris Sonchifolia Hance of Compositae and extracted as intravenous injection starting from full plants. The main pharmaceutical components of the intravenous injection are flavone and adenosine. The coexistence of phyto-flavone and -adenosine in Ixeris Sonchifolia Hance has shown a remarkable complementary effect for treatment of cardio-cerebral diseases, which has already been demonstrated by pharmacodynamics study and clinical application. The pharmacological effects of Ixeris Sonchifolia Hance injection consist in its efficacies of: 1. increasing coronary artery flow, lowering cardiovascular resistance, resisting myocardial infarction, enhancing collateral circulations, reducing myocardial oxygen consumption and improving cardio-microcirculations, which are useful for the treatment of coronary heart diseases, angina pectoris, chest distress and breath shortness, and myocardial infarction; 2. reducing platelet conglomeration, increasing the activity of fibrinolysin, inhibiting thrombosis, decreasing the viscosity of blood plasma and serum, increasing the electrophoresis velocity of erythrocytes, lowering cerebro-vascular resistance, increasing cerebral blood capacity and promoting restoration of neural function, which are useful in the treatment of cerebral infarction (cerebral thrombosis); 3. improving microcirculation disorders caused by bacteria, systemic microcirculation disorders caused by polymer dextran, and fundus microcirculation, and dilating fundus artery, which are useful in the treatment of fundus diseases, such as central retinitis, optic atrophy and retinitis pigmentosa. However, Ixeris Sonchifolia Hance injections exhibit poor stability and can only stand short-term storage. The contents of flavone and adenosine have been determined by High Performance Liquid chromatography (HPLC) and derivatization method, which show that adenosine content reduces from 15.0 \u03bcg/ml prior to finishing of the aqueous injection to 6.7 \u03bcg/ml thereafter, while for flavone, from 0.25 mg/ml to 0.169 mg/ml, with a lose of 55.3% and 32.4% of the prescribed, respectively. Six months later, adenosine content reduces to 6.5 \u03bcg/ml, while flavone content is 0.133 mg/ml. This can also be verified by subjecting Ixeris Sonchifolia Hance aqueous injection wherein flavone and adenosine contents are respectively 5.07 mg/ml and 24.37 \u03bcg/ml, to ten-day accelerated stress test carried out at 80\u00b0 C. in an oven, followed by content determination of both, in doing so, merely 3.16 mg/ml and 12.18 \u03bcg/ml are left for flavone and adenosine, respectively (Tab. 1). In view of the considerable lose of both active ingredients of flavone and adenosine during preparation and storage, the therapeutic effect of the injection is seriously impacted. Moreover, Ixeris Sonchifolia Hance injection exhibits a distinct change in color of solution. The absorbance measured at 400 nm using ultraviolet spectrophotometer during preparation of the aqueous injection from stock extract, ranges from 0.338 prior to sterilization to 0.423 thereafter, and further increases to 0.443 after 6 months with darkened color. It should be further noted that HPLC spectra of the injection exhibit poor similarity, that is to say, there exist major differences among HPLC spectra of different batches of aqueous injections yet formulated from the same stock extract. The stability behavior of same kind of commercial aqueous injections by random sampling (Batch No. 20000303, 20000503, 20010120) shows that: flavone contents are 0.27 mg/ml, 0.51 mg/ml and 0.30 mg/ml, while for adenosine, the contents are 1.11 \u03bcg/ml, 0.147 \u03bcg/ml and 0.00 \u03bcg/ml, respectively\nIn a conclusion, Ixeris Sonchifolia Hance aqueous injection is susceptible to various factors such as preparation, storage, and thus is difficult to be controlled in quality, which greatly expenses the therapeutic effect of Ixeris Sonchifolia Hance injection."}
-{"text": "The present invention relates to a process for the preparation of substituted dimethyl-(3-aryl-butyl)-amine compounds by means of homogeneous catalytic hydrogenation of dimethyl-(3-aryl-but-3-enyl)-amines.\nDimethyl-(3-aryl-butyl)-amine compounds have proved to be pharmaceutically active compounds exhibiting excellent analgesic activity and very good tolerability, see EP-A1-0 693 475. WO 2005/000788 A1 describes a process for the preparation of such compounds in which, in a first stage, substituted dimethyl-(3-aryl-but-3-enyl)-amine compounds are prepared by elimination of the tertiary hydroxyl group in substituted 4-dimethylamino-2-aryl-butan-2-ol compounds. These dimethyl-(3-aryl-but-3-enyl)-amine compounds are then hydrogenated in a second stage in the presence of metal catalysts. The heterogeneous hydrogenation proceeds in good yields with adequate activities. As expected, the stereoselectivity is not very pronounced. According to the examples in WO 2005/000788 A1, if two adjacent asymmetric C atoms are present diastereomeric ratios of from 2:1 to a maximum of 3:1 can be obtained for the trans diastereomer: cis diastereomer, that is to say always in favor of the trans diastereomer. The ratio is established automatically, essentially depends on the substrate, and can be influenced by the choice of reaction conditions to only a small extent."}
-{"text": "The area of miniaturized microfluidic technology, also known as \u201clab-on-a-chip\u201d or \u201cmicro total analysis systems\u201d, is rapidly expanding with the promise of revolutionizing chemical analysis and offering enabling tools and technologies for the life sciences. The majority of molecules and compounds of interest in the life sciences are in the liquid phase and typically analytical measurements are used to conduct quantitative and qualitative trace analysis of these analytes in solution. However, due to the nature and origin of the materials under analysis, sample amounts and volumes are historically in short supply and these amounts are constantly decreasing. Moreover, as seen in drug discovery and drug development, including pharmacokinetic and proteomics applications, the need for better analytical measurements that require smaller amounts and volumes of sample is growing. Inherently, microfluidics are a \u201cgood fit\u201d for the move to both smaller sample and volume requirements. In fact, a primary reason for miniaturization has been to enhance analytical performance of the device rather than reduce its physical size. Additionally, the ability to miniaturize the analysis with the use of microfluidics allows for integration of multiple separation techniques that enable parallel processing and also for the incorporation of several types of analytical measurements in a single device (sample handling, injections, 2D separations, reaction chambers etc.). Inherently, there are other benefits that accompany miniaturization, such as reduction in reagent and waste disposal, as well as, the reduction of the device footprint.\nThe first analytical miniaturized device was a gas chromatograph fabricated in silicon over 25 years ago. This device was designed for separating components in the gas phase. A decade later, components of a liquid chromatography column were fabricated on a silicon chip. Most early work related to miniaturization on silicon and the research focused on the fabrication of components such as microvalves, micropumps, and chemical sensors. None of the early systems implemented integrated electronics or electric fields for operation, but rather the silicon was used as a substrate for micromachining desired shapes/geometry.\nMost methods used in microfabrication were developed in the 1970's and 1980's in the silicon microprocessor industry. Typically, initial research developments were fabricated in silicon because of the extensive knowledge and tools available for silicon processing. This approach works satisfactorily for devices used for \u201cdry\u201d analyses, however, many microfluidic devices require the integration of on-chip electronics and/or the ability to apply electric fields to device. Because most applications in the life sciences involve samples contained in the liquid phase, the majority of micro analysis systems being developed/designed are for analysis of analytes contained in solution. The need for the ability to apply electric fields to the device becomes a serious issue when processing samples contained in the liquid phase on silicon substrates.\nA difficult scenario is encountered due to the opposing objectives of the micro total analysis system and the microprocessor technology used to make them. Typically, the microprocessor industry strives to keep microdevices \u201cdry\u201d and \u201cclean\u201d as liquids, moisture, and contaminates interfere with the device performance and operation. This highly contrasts the needs of micro total analysis systems where liquids and foreign substances (analytes, including salts) are deliberately introduced to the device. Again, this does not pose a problem for a silicon device that does not involve electronic or electrical field generation where only specific geometries are micromachined. However, major issues arise when electronics are incorporated in the devices and especially when potentials are applied for the generation of electric fields (semiconductor must be insulated for controlled electric field generation).\nThere is currently a move to perform chemical separations on-chip with the use of electric fields, for example applications such as CE, CEC, charged analyte manipulation, and charged solution manipulation. Because of the relatively strong fields needed for the separation process, research has moved to considering non-conductive substrates such as glass, quartz, and non-conducting polymers as opposed to the use of silicon as in the manufacture of semiconductors. This shift in materials is warranted because in order to form the electric fields, the substrates must be insulated in the desired areas.\nAlthough there are conventional techniques for insulating silicon substrates, the dielectric coatings currently available are designed for the electronics industry and operate under \u201cdry\u201d conditions. Much effort in the microprocessor industry has been expended on keeping devices dry or isolated from liquids. Additionally, the microprocessor industry has gone to great lengths to avoid contact of the electronic device with mobile ions such as salts due to the destructive nature they pose to dielectric coatings used to insulate the silicon used in microprocessors.\nMicrofluidic devices requiring the application of relatively high voltages and electric fields for the manipulation of liquids and samples are mainly fabricated on insulating substrates because of their insulating properties. Application of high voltages to liquids on insulators on conductor substrates often leads to shorting or drastically reduced performance and lifetime of desired electrical properties.\nAccordingly, the art needs dielectric coatings that do not degrade, but rather maintain their electrical properties when exposed to direct voltage application and high electric field strengths while in the presence of high humidity and/or direct liquid contact (wet). The art needs to overcome current coating technology limitations and provide appropriate solutions for microfluidic device applications. The art needs microfluidic devices that take advantage of the highly developed silicon processing techniques for silicon and other substrates including micromachining as well as electronic circuit integration and electric field definition. The art lacks the ability to incorporate microfluidics and electronics in the same substrate allowing for fully integrated systems."}
-{"text": "This invention relates to an evolutionary controlling system, and particularly to that for controlling characteristics of a subject in an evolutionary manner.\nHeretofore, when a control system or control characteristics of a subject, such as vehicles and electrical appliances, is designed, imaginary users are selected, and the users' preferences and their using conditions are taken into consideration. The characteristics of the subject are determined in such a way as to adapt the subject to users in as broad a range as possible.\nHowever, each individual user has a particular and unique personality, and thus, their preferences are diverse. Thus, there is a problem in that even if imaginary users are selected to develop and design a product for users by presuming the users' preference, it is impossible to satisfy all of the users of the product.\nIn order to solve the above problem, prior to purchase of a product, a prospective user is requested to determine whether or not the product is satisfactory to the user after checking the characteristics of the product in light of the user's preferences. However, it is troublesome for the user to check the characteristics of the product before purchase. Further, because a series of products are often operated or controlled by characteristics common in the products, although the design of the product is changed depending on the user's preferences, the user may not like other operational characteristics. Thus, although the design is appealing to some prospective users, the users may not purchase the product since the operational characteristics do not appeal to them. In other words, there is another problem in that the range of users is limited and depends on the operational characteristics.\nAn objective of the present invention is to provide an evolutionary control system to construct characteristics which can satisfy plural users."}
-{"text": "Aluminum metal has been produced for 90 years in the Hall cell by electrolysis of alumina in a molten cryolite salt electrolyte bath operating at temperatures in the range of 900.degree.-1000.degree. C. The reactivity of the molten cryolite, the need for excellent electrical conductivity, and cost considerations have limited the choice of materials for the electrodes and cell walls to the various allotropic forms of carbon.\nTypically the Hall cell is a shallow vessel, with the floor forming the cathode, the side walls a rammed coke-pitch mixture, and the anode a block suspended in the bath at an anode-cathode separation of a few centimeters. The anode is typically formed from a pitch-calcined petroleum coke blend, prebaked to form a monolithic block of amorphous carbon. The cathode is typically formed from a pre-baked pitch-calcined anthracite or coke blend, with cast-in-place iron over steel bar electrical conductors in grooves in the bottom side of the cathode.\nDuring operation of the Hall cell, only about 25% of the electricity consumed is used for the actual reduction of alumina to aluminum, with approximately 40% of the current consumed by the voltage drop caused by the resistance of the bath. The anode-cathode spacing is usually about 4-5 cm., and attempts to lower this distance result in an electrical discharge from the cathode to the anode through aluminum droplets.\nThe molten aluminum is present as a pad in the cell, but is not a quiescent pool due to the factors of preferential wetting of the carbon cathode surface by the cryolite melt in relation to the molten aluminum, causing the aluminum to form droplets, and the erratic movements of the molten aluminum from the strong electromagnetic forces generated by the high current density.\nThe wetting of a solid surface in contact with two immiscible liquids is a function of the surface free energy of the three surfaces, in which the carbon cathode is a low energy surface and consequently is not readily wet by the liquid aluminum. The angle of a droplet of aluminum at the cryolite-aluminum-carbon junction is governed by the relationship ##EQU1## where .alpha..sub.12, .alpha..sub.13, and .alpha..sub.23 are the surface free energies at the aluminum carbon, cryolite-carbon, and cryolite-aluminum boundaries, respectively.\nIf the cathode were a high energy surface, such as would occur if it were a ceramic instead of carbon, it would have a higher contact angle and better wettability with the liquid aluminum. This in turn would tend to smooth out the surface of the liquid aluminum pool and lessen the possibility of interelectrode discharge allowing the anode-cathode distance to be lowered and the thermodynamic efficiency of the cell improved, by decreasing the voltage drop through the bath.\nTypically, amorphous carbon is a low energy surface, but also is quite durable, lasting for several years duration as a cathode, and relatively inexpensive. However, a cathode or a TiB.sub.2 stud as a component of the cathode which has better wettability and would permit closer anode-cathode spacing could improve the thermodynamic efficiency and be very cost-effective.\nSeveral workers in the field have developed refractory high free energy material cathodes. U.S. Pat. No. 2,915,442, Lewis, Dec. 1, 1959, claims a process for production of aluminum using a cathode consisting of the borides, carbides, and nitrides of Ti, Zr, V, Ta, Nb, and Hf. U.S. Pat. No. 3,028,324, Ransley, Apr. 3, 1962, claims a method of producing aluminum using a mixture of TiC and TiB.sub.2 as the cathode. U.S. Pat. No. 3,151,054, Lewis, Sept. 29, 1964, claims a Hall cell cathode conducting element consisting of one of the carbides and borides of Ti, Zr, Ta and Nb. U.S. Pat. No. 3,156,639, Kibby, Nov. 10, 1964, claims a cathode for a Hall cell with a cap of refractory hard metal and discloses TiB.sub.2 as the material of construction. U.S. Pat. No. 3,314,876, Ransley, Apr. 18, 1967, discloses the use of TiB.sub.2 for use in Hall cell electrodes. The raw materials must be of high purity particularly in regard to oxygen content, Col. 1, line 73-Col. 2, line 29; Col. 4, lines 39-50, Col. 8, lines 1-24. U.S. Pat. No. 3,400,061, Lewis, Sept. 3, 1968 discloses a cathode comprising a refractory hard metal and carbon, which may be formed in a one-step reaction during calcination. U.S. Pat. No. 4,071,420, Foster, Jan. 31, 1978, discloses a cell for the electrolysis of a metal component in a molten electrolyte using a cathode with refractory hard metal TiB.sub.2 tubular elements protruding into the electrolyte. The protruding elements enhance electrical conductivity and form a partial barrier to the mechanical agitation caused by magnetic effects."}
-{"text": "The present invention relates to progesterone receptor modulators.\nIntracellular receptors (IR) form a class of structurally related gene regulators known as \u201cligand dependent transcription factors\u201d. The steroid receptor family is a subset of the IR family, including progesterone receptor (PR), estrogen receptor (ER), androgen receptor (AR), glucocorticoid receptor (GR), and mineralocorticoid receptor (MR). A compound that binds to an IR and mimics the action of the natural hormone is termed an agonist, whilst a compound that inhibits the effect of the hormone is an antagonist.\nThe natural hormone, or ligand, for the PR is the steroid progesterone, but synthetic compounds, such as medroxyprogesterone acetate or levonorgestrel, have been made which also serve as ligands. Once a ligand is present in the fluid surrounding a cell, it passes through the membrane via passive diffusion, and binds to the IR to create a receptor/ligand complex. This complex binds to specific gene promoters present in the cell's DNA. Once bound to the DNA the complex modulates the production of mRNA and protein encoded by that gene.\nPR agonists (natural and synthetic) are known to play an important role in the health of women. PR agonists are used in birth control formulations, typically in the presence of ER agonists, alternatively they may be used in conjunction with PR antagonists. ER agonists are used to treat the symptoms of menopause, but have been associated with a proliferative effect on the uterus that can lead to an increased risk of uterine cancers. Co-administration of a PR agonist reduces or ablates that risk.\nU.S. Pat. No. 6,407,101, which is hereby incorporated by reference, describes the preparation of cyclocarbamate derivatives, which are useful as progesterone receptor modulators. These cyclocarbamate derivatives, including, e.g., 5-(4,4-dimethyl-2-thioxo-1,4-dihydro-2H-benzoxazin-6-yl)-1-methyl-1H-2-cyano-pyrrole, are prepared by thionation of the corresponding benzoxazin-2-one (Scheme 1).\n\nWhat is needed in the art are alternate compounds that are effective as progesterone receptor modulators."}
-{"text": "Actuators typically are mechanical devices that are used for moving or controlling a mechanism, system or the like and typically convert energy into some type of motion. Examples of actuators can be found in any number of applications encountered in everyday life including automotive, aviation, construction, farming, factories, robots, health care and prosthetics, among other areas.\nMobile robotics and advanced prosthetics will likely play important roles in the future of the human race. Actuators frequently are used in these applications that enable movement of a robot or user arm or other appendage or item as desired.\nMost existing mobile robots and advanced prosthetics, however, lack the strength and speed necessary to be effective. This is because they suffer from poor specific power (strength\u00d7speed/weight) which determines how quickly work can be done compared to another actuator of the same weight.\nFor example, if such devices are capable of lifting significant weight, they must do so very slowly, which inhibits their adoption for most applications. On the other hand, devices that can move more quickly are just not capable of handling anything more than the smallest weight."}
-{"text": "The present invention relates generally to a connector and, more particularly, to a unique tool for removing contacts from a connector body.\nWhile the present invention will be described specifically in connection with a tool for removing contacts from an electrical connector, it will be appreciated that the tool could also be utilized for removing ferrules from fiber optic connectors.\nIn U.S. Pat. No. 3,110,093 to Johnson, there is disclosed an electrical connector of the type wherein one or more contacts, each connected to a wire, are inserted from the rear into a contact receiving bore in a connector body or insulator after the connector has been otherwise completely fabricated or assembled. The insulator includes a retaining clip or other locking means between the individual contacts and their respective bore walls for retaining the contacts in their operative position in the insulator. The Johnson patent also discloses a suitable tool for insertion into clearance space between the contact and the bore wall from the rear of the insulator to disengage the locking device and thereby permit manual withdrawal of the contact from the rear of the insulator by pulling on the wire attached to the contact. While this arrangement works well when a wire is present, the tool is ineffective fr an unwired contact in that the wire receiving portion of the contact normally does not extend beyond the rear of the insulator to permit the contact to be gripped for withdrawal.\nU.S. Pat. No. 3,380,141 Rofer discloses a tool for removing an unwired contact which, when inserted into the contact bore from the rear of the insulator, releases the locking means therein and frictionally engages the rear of the contact so that it may be withdrawn from the bore.\nFrench Pat. No. 2,240,600 discloses a front release contact retention assembly, in contrast to the rear release arrangments disclosed in the aforementioned Johnson and Rofer patents, wherein a contact extraction tool is inserted through the contact bore from the front of the insulator to release the contact locking means therein. The tip of the tool is longitudinally slotted and dimensioned so that it will frictionally engage the contact. Thus, as in the Rofer arrangement, the contact locking means is released while the tool frictionally engages the contact, thereby allowing the contact to be removed from the insulator. If manufacturing tolerances are not closely held in the forming of the prior art tools, it will be appreciated that adequate frictional engagement of the tool with the contact may not occur with the result that the tip of the tool may slip off the contact when the tip is withdrawn from the contact bore. Furthermore, even if the tool initially frictionally engages the contact with adequate force to allow the contact to be withdrawn from the insulator, continued use of the tool can result in the tip of the tool wearing so that it no longer firmly engages the contact.\nIt is the object of the present invention to provide a contact extraction tool which embodies manually operable gripping jaws which will assure that a firm gripping force can be applied to the contact so that the contact will be pulled out of the insulator when the tip of the tool is withdrawn from the contact bore."}
-{"text": "The polyphenylene ether resins are a family of engineering thermoplastics that are well known to the polymer art. These polymers may be made by a variety of catalytic and non-catalytic processes from the corresponding phenols or reactive derivatives thereof. By way of illustration, certain of the polyphenylene ethers are disclosed in Hay, U.S. Pat. Nos. 3,306,874 and 3,306,875, and in Stamatoff, U.S. Pat. Nos. 3,257,357 and 3,257,358. In the Hay patents, the polyphenylene ethers are prepared by an oxidative coupling reaction comprising passing an oxygen-containing gas through a reaction solution of a phenol and a metal-amine complex catalyst. Other disclosures relating to processes for preparing polyphenylene ether resins, including graft copolymers of polyphenylene ethers with styrene type compounds, are found in Fox, U.S. Pat. No. 3,356,761; Sumitomo, U.K. Pat. No. 1,291,609; Bussink et al. U.S. Pat. No. 3,337,499; Blanchard et al. U.S. Pat. No. 3,219,626; Laakso et al. U.S. Pat. No. 3,342,892; Borman, U.S. Pat. No. 3,344,166; Hori et al. U.S. Pat. No. 3,384,619; Faurote et al. U.S. Pat. No. 3,440,217; and disclosures relating to metal based catalysts which do not include amines, are known from patents such as Wieden et al., U.S. Pat. No. 3,442,885 (copper-amidines); Nakashio et al., U.S. Pat. No. 3,573,257 (metal-alcoholate or -phenolate); Kobayashi et al. U.S. Pat. No. 3,455,880 (cobalt chelates); and the like. In the Stamatoff patents, the polyphenylene ethers are produced by reacting the corresponding phenolate ion with an initiator, such as peroxy acid salt, an acid peroxide, a hypohalite, and the like, in the presence of a complexing agent. Disclosures relating to non-catalytic processes, such as oxidation with lead dioxide, silver oxide, etc. are described in Price et al., U.S. Pat. No. 3,382,212. Cizek, U.S. Pat. No. 3,383,435 discloses polyphenylene ether-styrene resin compositions. All of the above-mentioned disclosures are incorporated by reference.\nThe term \"polystyrene resin\" includes polymers and copolymers of styrene, alpha methyl styrene, chlorostyrene, and the like.\nThe term \"EPDM\" includes rubbery interpolymers of a mixture of mono-olefins and a polyene. Preferred types are those rubbery interpolymers of ethylene, an alpha-olefin, and a polyene. Rubbery interpolymers of ethylene, propylene, and a polyene are especially preferred.\nThe term \"EPDM-silicone rubber\" includes graft copolymers of EPDM with silicone rubber. An example of an EPDM-silicone rubber employed in the present invention is that commercially manufactured and sold by Shinetsu Chemical Industry, under the name SEP-172U.\nIn the prior art, rubber-modified styrene resins have been admixed with polyphenylene ether resins to form compositions that have modified properties. The Cizek patent, U.S. Pat. No. 3,383,435, discloses rubber-modified styrene resin-polyphenylene ether resin compositions wherein the rubber component is of the unsaturated type such as polymers and copolymers of butadiene. The physical properties of these compositions are such that it appears that many of the properties of the styrene resins have been upgraded, while the moldability of the polyphenylene ethers are improved.\nNakashio et al. U.S. Pat. No. 3,658,945 discloses that from 0.5 to 15% by weight of an EPDM-modified styrene resin may be used to upgrade the impact strength of polyethylene ether resins. In Cooper et al., U.S. Pat. No. 3,943,191 it is disclosed that when the highly unsaturated rubber used in compositions of the type disclosed by Cizek, is replaced with EPDM rubber that has a low degree of residual unsaturation, the thermal oxidative stability and color stability are improved. The EPDM rubber in the Cooper et al. compositions is comprised substantially of particles in the range of 3-8 microns.\nThe impact strength of the Cooper et al. compositions is superior to that of a polypropylene ether resin alone or that of similar compositions comprised of unmodified polystyrene; however, the impact strength of the Cooper et al. compositions is inferior to that of similar compositions comprised of polystyrene modified with polybutadiene rubber, such as a composition known as FG-834 available from Foster-Grant Co. As is disclosed in U.S. Pat. No. 3,981,841, the impact strength of the Cooper et al. compositions can be improved by incorporating therein impact modifiers such as an emulsion-grafted EPDM polystyrene copolymer. U.S. Pat. No. 4,152,316 incorporated herein by reference, discloses that a composition of a polyphenylene ether resin and an alkenyl aromatic resin modified with an EPDM rubber comprised of particles having a median diameter less than about two microns, preferably about 0.5 to 1.5 microns, is a very useful thermoplastic molding material having good thermal oxidative stability and good room temperature impact strength but inferior low temperature impact strength.\nIn U.S. Pat. No. 4,102,850 it is disclosed that the addition of small amounts of mineral oil to the polymerizing mixture of styrene and EPDM rubber produces EPDM-modified polystyrene which yields blends with polyphenylene oxide having significantly better low-temperature impact strength than blends made from EPDM-polystyrene made without the mineral oil.\nIn U.S. Pat. No. 3,737,479 it is disclosed that the addition of a silicone such as a polyorganosiloxane, which is fluid, to polyphenylene oxide or to polyphenylene oxide-polystyrene blends improves Gardner impact strength but not Izod impact strength (see Col. 1, lines 54-57 of U.S. Pat. No. 3,737,479).\nIt has now been found when small amounts of a graft copolymer of EPDM with silicone rubber is blended with a polyphenylene oxide resin, and optionally with a polystyrene resin, blends can be obtained which are substantially improved in ductility and Izod impact strength both at room temperature and at low temperatures."}
-{"text": "This invention is in the field of terminally guided, anit-armor missiles. Heretofore known terminal guidance missile systems have used proportional navigation with limited trajectory shaping for high accuracy against moving targets. This limited use of trajectory shaping results in either a flat approach trajectory which has reduced warhead penetration or a lofted or ballistic like trajectory. The ballistic like trajectory is often unable to perform well when low cloud cover condition exists. The including of a gimbal angle regulator and boost/cruise trajectory shaping allows low altitude ground or air launch, climb to cruise altitude under a low cloud cover, and then dive onto the target thereby achieving a high probability of penetrating the armor of the target. This invention is not limited to tracker systems that must acquire the target prior to launch such as an infrared imaging seeker, but works equally well with systems that can acquire the target after launch such as laser semiactive systems. The use of conventional known guidance schemes cannot accomplish the high probability of accurate hit concurrent with control of the dive angle for maximum warhead performance. Any ballistic like trajectory for terminally guided missiles must reacquire after descending through the cloud cover adversely affecting the probability of hit and limiting the controllability of the impact attitude."}
-{"text": "This invention is concerned with an adapter to enable a dry cell battery to be used in an application in which it might be charging or discharging for which it is electrically acceptable but for which it is physically too small.\nBatteries in flash lights, portable radios or so called chargers and other equipment are usually supported in a tube or clips or other mechanical contrivances so that they are securely held in electrical contact with the terminals of the circuit or element which they are to supply or from which they are to be supplied (charged). The support means may hold a single battery between the terminals or a plurality of batteries in end to end relationhip may be so held but almost always the support means are arranged to receive a battery of particular transverse cross sectional dimensions.\nAs is well known, there are generally available several differently dimensioned batteries which have similar electrical discharge and charging characteristics. It frequently occurs than when a physically large battery is needed there is available only a battery of a smaller size but one which is electrically acceptable.\nFurther, as batteries are reduced in size with advancing technology, it becomes technically and economically desirable to use the new and smaller batteries in equipment designed specifically to receive older and larger batteries. For example, a nickel/cadmium AA battery has similar electrical discharge characteristics to a zinc/carbon D battery but of course as it is rechargeable it is desirable to be able to use such a battery in place of the larger zinc/carbon battery.\nMajor differences between existing large and small batteries is their working lives. Many pieces of equipment such as, for example, flash lights are designed to receive large batteries that fit them for long and continuous use but in practice those pieces of equipment may be used for only short periods with long intervals of non-use between those periods. During those periods of non-use the battery deteriorates or, because of faulty terminals in the equipment, discharges. In many cases these long periods of non-use can be foreseen but because of the physical design characteristics of the pieces of equipment one is obliged to use a large and expensive battery knowing that were its use physically possible a small, cheaper battery having a shorter working like would be quite adequate."}
-{"text": "Fuel injected engines employ fuel injectors, each of which delivers a metered quantity of fuel to an associated engine cylinder during each engine cycle. Prior fuel injectors were of the mechanically or hydraulically actuated type with either mechanical or hydraulic control of fuel delivery. More recently, electronically controlled fuel injectors have been developed. In the case of an electronic unit injector, fuel is supplied to the injector by a transfer pump. The injector includes a plunger which is movable by a cam-driven rocker arm to compress the fuel delivered by the transfer pump to a high pressure. An electrically operated mechanism either carried outside the injector body or disposed within the injector proper is then actuated to cause fuel delivery fuel to the associated engine cylinder.\nIn prior fuel injector designs, high pressure fuel is conducted through passages which are located outside of a central recess containing a solenoid which operates a valving mechanism. The passages are located close to the outer surface of the fuel injector and are formed by drilling intersecting holes. After drilling, portions of some of the holes must be filled with plugs. These passages and plugs are subjected to very high fluid pressures, thus requiring careful design, thus increasing complexity and cost.\nIn addition to the foregoing, because the high pressure passages are located outside of the solenoid, the size of the solenoid is necessarily limited, thereby limiting the available solenoid force.\nStill further, a prior type of fuel injector utilizes a direct operated check valve, which includes upper and lower valve seats which must be precisely aligned for proper operation. Manufacturing and assembly tolerances must, therefore, be kept tight, further increasing cost."}
-{"text": "Hydroxyl terminated polyethers find uses as raw materials of polyurethane and polyester resins. Others uses include cosmetics, plasticizers, surfactants and raw materials of these products. Conventional polyethers are terminated with primary or secondary alcoholic groups, the presence of which normally contributes to their desirablity for intended uses. It has been experienced in certain cases that the primary or secondary alcohol function is too active and therefore the presence of which is not desirable. For example, polyethers having too active an alcoholic function can lead to a premature reaction when reacted with polyisocyanates. Attempts have been made to remove the hydroxyl function by blocking with a group such as methyl or acyl, or replacing the hydroxyl group with chlorine. Obviously, polyethers having blocked or chlorinated terminals are not suitable for uses in which the presence of a hydroxyl function is essential, e.g. for the production of polyurethane and polyester resins."}
-{"text": "The present invention relates generally to a computer implemented method for selecting responsive actions. More specifically, the present invention relates to referencing earlier attempted actions, if sufficiently current, and responding with at least one alternate action.\nSystems management software installed on data processing systems are used to detect and catalog events. An event is a signal that indicates a condition of a resource of a data processing system. An event can be, for example, a condition related to a rotation rate of a fan. Accordingly, the event can be a signal generated to identify both the resource and the condition of the resource. Some events can be asynchronously generated from an error condition. In addition, multiple events may be aggregated together to form an event summary. An event summary is an event that describes at least two other events. For example, an event summary can be a count of the number of times a fan speed exceeded a threshold in the previous hour. A resource is a physical data processing hardware or virtual resource, backed by physical data processing hardware. A resource can correspond to a system, subsystem, software program, or device. A resource can have attributes such as a temperature. An event, on the other hand, may be a request such as a hypertext protocol request at a server, or a report that a processor on that server has met or exceeded a temperature. Accordingly, the use of events can be helpful to collect information about a data processing system's status.\nOne or more events can be indicative of an underlying problem or anomaly. Accordingly, it is helpful that system administrators pay attention to some events under some conditions."}
-{"text": "1. Field of the Invention\nThis invention relates generally to memory design. Particularly, this invention relates to a new cell for implementation in the Static Random-Access Memory (SRAM) thus taking advantage of the base current reversal phenomenon in a Complementary Metal Oxide Semiconductor (CMOS) for the construction of compatible high gain gated lateral Bipolar Junction Transistor (BJT).\n2. Description of the Prior Art\nThe conventional CMOS SRAM cell essentially consists of a pair of cross-coupled inverters as the storage flip-flop or latch, and a pair of pass transistors as the access devices for data transfer into and out of the cell. (A large number of different cell configurations are cited in the literature, e.g.,B. Prince, Semiconductor Memories--A Handbook of Design, Manufacture, and Application, 2nd ed., New York: John Wiley & Sons, Inc., 1991). Thus, a total of six Metal Oxide Semiconductor Field Effect Transistors (MOSFETs) or four MOSFETs plus two very high resistance load devices are required for implementing a conventional CMOS SRAM cell. However, to achieve high packing density, it is the usual practice to reduce the number of the devices needed for realizing a CMOS SRAM cell or the number of the devices for performing the Write/Read operation. Especially for the case of very high resistance load devices, increased process complexity, extra masks, and high fabrication cost are required for forming the undoped polysilicon layers or the Thin Film Transistor (TFT) on the oxide and thus saving the chip area; however, the corresponding product yield is not high. Therefore, more efforts are needed to further reduce the areas occupied by the chip while improving the production yield.\nThe structures of the conventional SRAM are shown in FIGS. 1 and 2. FIG. 1 shows a circuit schematic of a conventional SRAM cell configuration. The cell comprises a pair of cross-coupled inverters, p-MOSFET 11 and n-MOSFET 13, and p-MOSFET 16 and n-MOSFET 14, as the storage flip-flop or latch. In each inverter the gates of p-MOSFET and n-MOSFET are tied together and connected to the output of another inverter. The output of each inverter are the drains of p-MOSFET and n-MOSFET which are tied together. The conventional cell employs a pair of pass transistors, n-MOSFETs 12 and 15, as the access devices for data transfer into and out of the cell. Two column lines DATA and DATA and two row select lines W and W are depicted. This conventional cell requires 6 MOSFETs.\nFIG. 2 shows a circuit schematic of a conventional CMOS cell using undoped polysilicon layer or thin-film transistor for providing very high resistance loads 20. FIG. 2 is very similar to FIG. 1 except for the type of the load. This cell requires complex processes and the resulting yield is low. This conventional cell requires four MOSFETs formed on the silicon plus two undoped polysilicon layers or TFTs formed on the oxide.\nA Bipolar-CMOS (BICMOS) process has recently been introduced to achieve the above-mentioned object (refer to the paper: K. Sakui, et al. \"A new static memory cell based on reverse base current (RBC) effect of bipolar transistor,\" IEEE IEDM Tech. Dig., pp. 44-47, December 1988). In this BICMOS process, only two devices are needed for a SRAM cell: one vertical bipolar transistor and one MOSFET as a pass device. However, extra processing steps and increased masks are required along with special deep isolation techniques, resulting in high fabrication cost and process complexity. Also, the yield of the SRAM products utilizing such complex BICMOS process is usually low compared with the existing standard CMOS process. Therefore, the SRAM products fabricated in a complex BICMOS process cannot provide sufficient competition to the conventional CMOS SRAM products.\nVery recently, a new phenomenon of base current reversal has been observed in a CMOS compatible high gain gated lateral bipolar transistor (refer to the paper: Tzuen-Hsi Huang and Ming-Jer Chen, \"Base current reversal phenomenon in a CMOS compatible high gain n-p-n gated lateral bipolar transistor,\" IEEE Trans. Electron Devices, Vol. 42, No. 2, pp.321-327, February 1995). This new phenomenon has been extensively investigated and has been found to have many applications. This invention employs this new phenomenon in the existing standard CMOS process for realizing a new SRAM cell constructed by only two MOSFETs. This new SRAM also features only one-sided peripheral circuitry for Read/Write action. Therefore, this invention can efficiently save the chip area with high yield since it is fully compatible with the existing low-cost standard CMOS process."}
-{"text": "Many types of electrical equipment contain IC (integrated circuit) devices which are vulnerable to damage from high voltage transients.\nIn a television receiver the anode of the image producing kinescope is typically biased at a high potential, e.g., 25,000 volts. High voltage transients may be produced when the high voltage anode of the kinescope is rapidly discharged to points at lower potentials. Such high voltage transients have positive and negative peaks often in excess of 100 volts and may last several microseconds. High voltage transients may also be produced when electrostatic charges are discharged as a user contacts the controls of the television receiver. High voltage transients may be coupled to the terminals of IC's employed in the television receiver for video and audio signal processing. Accordingly, these IC's may be damaged by high voltage transients.\nCopending patent application entitled \"Protection Circuit for Integrated Circuit Devices,\" Ser. No. 212,534, filed in the name of the present inventor Dec. 3, 1980 and assigned to RCA Corporation, discloses a protection device for protecting an IC from negative voltage transients that exceed the most negative power supply potential applied to the IC. Another copending patent application entitled \"Protection Circuit for Integrated Circuit Devices,\" Ser. No. 230,357 filed also in the name of the present inventor Jan. 30, 1981, and also assigned to RCA Corporation, discloses a protection circuit which is triggered at one forward biased diode voltage drop above the positive supply voltage for protecting an IC from positive voltage transients that exceed the most positive power supply potential applied to the IC.\nHowever, in a television receiver, particular signals applied to an IC may have positive voltage excursions which in normal operation exceed the positive supply potential. For example, a typical television horizontal/vertical regulator IC requires a feedback connection from the kinescope deflection coils to one of its input terminals. While the power supply for the IC is typically +10 volts, the peak feedback voltage from the deflection coils is typically +27 volts. Therefore, it is desirable to provide a positive transient protection circuit for such IC's that permits normal signal voltages to exceed the power supply potential without activating such protection circuit and nevertheless protects the IC from excessively large transients."}
-{"text": "A common construction of a header for an agricultural implement includes either a cutter bar with a pickup reel and/or a rotating tine pickup belt for feeding material to an auger which collects the material into a rear discharge of the header. In certain types of crops, material can be thrown upward onto the frame of the header above the rear discharge opening. The material can then obstruct the view of the operator of the implement, or become entangled in the implement, or be wasted by being thrown onto the ground instead of being collected in the usual manner of the implement."}
-{"text": "Not Applicable\nNot Applicable\nThe present invention relates to fan drive assemblies of the type including a cooling fan and a fan drive, and more particularly, to such fan drive assemblies wherein the fan drive is of the type in which heat is generated as a result of the transmission of torque within the fan drive, and the ability of the fan drive to dissipate such generated heat represents a limiting factor on the torque transmitting capability of the fan drive assembly.\nAlthough the present invention may be used with various types and configurations of torque transmitting fan drives, it is especially adapted for use with fan drive assemblies of the type including a viscous fluid coupling device as the fan drive, and will be described in connection therewith.\nFan drive assemblies of the type which may benefit from the use of the present invention have found several uses, one of the most common of which is in connection with cooling the radiator of a vehicle engine. As is well known to those skilled in the art, the fan drive of the typical fan drive assembly comprises a viscous fluid coupling device, so named because the coupling utilizes a high viscosity fluid to transmit torque, by means of viscous shear drag, from an input coupling member (clutch) to an output coupling member (housing), with the cooling fan being bolted, or otherwise suitably attached, to the output coupling member.\nThe present invention is especially advantageous when used on a relatively high horsepower fan drive assembly, i.e., one which is capable of transmitting somewhere in the range of about two to about twelve horsepower from the fan drive to the cooling fan. Typically, such high horsepower fan drives include an output coupling member of the type which comprises a cast aluminum body and a cast aluminum cover. The input coupling member is typically also made as a cast aluminum member, and cooperates with the body and/or the cover to define a plurality of interdigitated lands and grooves which define the viscous shear space. When the shear space is filled with viscous fluid, typically a silicon fluid, torque is transmitted from the input coupling member to the output coupling assembly, in response to the rotation of the input coupling member.\nDuring such torque transmission, substantial heat is generated as a result of the shearing of the viscous fluid between the lands and grooves. The amount of heat generated is generally proportional to the xe2x80x9cslip speedxe2x80x9d of the fan drive, i.e., the difference between the speed of the input and the speed of the output. It is generally well understood by those skilled in the art that the ability to transmit torque is limited by the ability of the device to dissipate the heat generated. For example, in a viscous fan drive, if the temperature of the viscous fluid exceeds a certain maximum temperature, the result will be a deterioration of the viscous properties of the fluid, resulting in a gradual loss of the torque transmitting capability of the fluid.\nIn the fan drive art, it has been conventional for the design and development of a particular cooling fan to occur generally independently of the design and development of the viscous fan drive with which the fan is to be utilized. In other words, the fan is designed to provide the desired operating parameters (e.g., torque, air flow, etc.), thus determining the blade configuration and spacing, and then the mounting portion of the fan (the xe2x80x9cspiderxe2x80x9d) is designed or merely modified to adapt to the configuration of the particular fan drive mounting arrangement (e.g., mounting pads or bosses, disposed at a particular diameter from the axis of the fan drive).\nWhat has not been conventional in the fan drive art is to design the cooling fan and the fan drive as a xe2x80x9cpackagexe2x80x9d, with the goal of maximizing the heat dissipation of the overall fan drive assembly. As a result, it would appear that, at the time of the present invention, there is no commercially available fan drive assembly which achieves nearly its optimum, potential heat dissipation (heat rejection). As a further result, practically every fan drive assembly in commercial use is larger and more expensive than is actually necessary, in order to achieve a particular, desired flow of cooling air through the radiator.\nAlthough the present invention is not limited to a fan drive assembly in which the fan is mounted to the rearward side of the housing (body), rather than being mounted to the cover, the invention is especially advantageous in such an arrangement, and will be described in connection therewith. A typical xe2x80x9crear mountxe2x80x9d fan drive is illustrated and described in U.S. Pat. No. 4,384,824, in which the body member includes four mounting bosses located radially inward of the body cooling fins. As a result, the fan spider interferes with the radial flow of cooling air through the body cooling fins, thus reducing the heat dissipation capability of the particular fan drive assembly.\nAccordingly, it is an object of the present invention to provide an improved fan drive assembly in which the cooling fan and the fluid coupling device driving the fan are designed such that the overall assembly approaches the optimum, potential heat dissipation.\nIt is a more specific object of the present invention to provide an improved fan drive assembly in which the cooling fan is mounted to the body (housing) of the fluid coupling device in a manner which substantially improves the flow of air through the housing cooling fins.\nIt is a related object of the present invention to provide an improved fan drive assembly which accomplishes the above-stated objects, and in which the cooling fan is configured to further improve the flow of air through the housing cooling fins.\nIt is another object of the present invention to provide an improved fan drive assembly which accomplishes the above-stated objects, and in which the cover cooling fins are configured to improve the flow of air through the cover cooling fins.\nThe above and other objects of the invention are accomplished by the provision of a fan drive assembly of the type comprising a cooling fan attached to a fluid coupling device, the cooling fan comprising a fan hub, a spider portion, and a plurality of fan blades extending radially from the fan hub. The fluid coupling device comprises a first rotatable coupling assembly including a body member having a rearward surface, and a cover member cooperating with the body member to define a fluid chamber therebetween, a second rotatable coupling member being disposed in the fluid chamber for rotation relative to the first coupling assembly. The first coupling assembly and the second coupling member cooperate to define a viscous shear chamber therebetween, whereby torque may be transmitted from the second coupling member to the first a coupling assembly in response to the presence of viscous fluid in the shear chamber. The body member includes a plurality of cooling fins and a plurality of mounting portions, the spider portion being attached to the mounting portions and defining a pilot diameter.\nThe improved fan drive assembly is characterized by the body member including a plurality of mounting portions, each of which is disposed immediately adjacent an outer periphery of the body member. Each of the mounting portions defines a machining chucking surface, and a spider mounting surface on a rearward face thereof. The spider mounting surface includes a pilot surface in engagement with the pilot diameter of the spider portion. The plurality of cooling fins covers substantially all of the rearward surface of the body member not covered by the mounting portions.\nIn accordance with another aspect of the present invention, the fan hub and each of the plurality of fan blades cooperate to define a rearward axially extending air dam portion operable to restrict localized radial air flow."}
-{"text": "In thermal dye transfer printing, an image is formed on a receptor sheet by selectively transferring an image forming material to the receptor sheet from a dye donor sheet. Material to be transferred from the dye donor sheet is selected by a thermal printhead, which consists of small, electrically heated elements. These elements transfer image-forming material from the dye donor sheet to areas of the dye receptor sheet in an imagewise manner.\nThere are three broad classes of thermal transfer systems that are known, (1) chemical reaction systems, (2) thermal mass transfer systems, and (3) thermal dye transfer systems.\nIn chemical reaction systems, the image is formed upon the receptor as a result of the imagewise transfer of some chemical reactant from the donor sheet. An example is the transfer of a mobile molecule, such as phenol, to the receptor sheet, which bears a leuco compound thereon. The phenol is transferred by being volatilized by the heat from the thermal print head, and, upon reaching the receptor sheet, reacts with the leuco compound to convert it from the colorless to the colored form. Alternately, the phenol can be on the receptor sheet and the leuco compound can be on the donor sheet.\nIn thermal mass systems, no color forming chemical reaction takes place. Instead, the image is formed simply by the transfer of a mass of material containing colorant therein, such as pigment-filled polymer coatings.\nIn thermal dye transfer systems, a dye donor sheet is used in combination with a dye receptor sheet wherein, with the application of heat, a dye is transferred onto the receptor sheet at a controlled amount to obtain a dye image having gradation like in a photograph.\nEach system has its own advantages and disadvantages for the particular application of thermal printing. Various problems have been encountered with each proposed system. For thermal dye transfer, dye release layers have been proposed to enable efficient transfer of the dye layer from the dye donor sheet. Also, various dye-permeable release layers on the dye receptor layer have been proposed. The dye-permeable release layer is coated over the dye receptor layer, and is formulated to prevent sticking between the donor layer and the receptor layer during the transfer of the dye across the binder membranes. The release layer must also be formulated to allow effective transfer of the dye through the release layer. In general, many of these problems have been related to a specific resin used in the composition of the dye donor or dye receptor layers.\nSelection of the functional resin systems for the dye donor and the receptor sheet layers has been the topic of concern for many proposed dye systems. In consideration of the above mentioned requirements, efficient dye transfer and sticking between the dye donor layer and the dye receptor layer during transfer, a good, functional resin system to eliminate some of these problems is needed. Interestingly, a unique resin system has been found that provides an efficient working dye donor sheet and dye receptor sheet. The resin system comprises a chlorinated polyvinyl chloride (CPVC).\nU.S. Pat. No. 3,584,576 describes a heat sensitive stencil sheet comprising a film adhered to a porous thin fibrous sheet. The stencil sheet is perforated by exposure to infrared rays. The film consists essentially of at least 75% by weight of a chlorinated polyvinyl chloride resin, the balance being a polyvinyl chloride resin. A colorant may be present in the film. Upon being heated by infrared radiation, the film melts and forms perforations. The pores in the remaining fibrous sheet enable stencilling to be done through the perforations and the sheet.\nThere are noticable difference between the above mentioned prior art use of CPVC in a thermally sensitive stencil applications and in the present invention. The prior art uses CPVC merely as a resinous binder, with or without other resinous binders. It is, in particular, used not as a receptor layer for a thermal dye transfer sheet, but as the thermoplastic binder for a thermal stencil sheet. The novel use of the CPVC resin in the thermal dye transfer printing of the present invention has been found to give surprisingly new and unique properties for use as the primary resinous thermal plastic binder in both a dye donor sheet, and a dye receptor sheet. Typically, commercially available dye donor sheets and dye receptor sheets are comprised of chemically different binders with different functionals."}
-{"text": "The invention relates to a method for the manufacture of a wire electrode for the spark-erosion process and a wire electrode for this method, which consists of a core of steel, an intermediate layer of copper or a high copper-containing alloy, and an outer layer with at least 40% zinc or a core of steel and a zinc alloy outer layer having a zinc content of from 40-60%.\nWire electrodes, which are designed with multiple layers, which have a core of steel, an intermediate layer of copper arranged around the core, and an outer zinc-containing layer, are known, for example, from the DE-PS 29 06 245. All known wire electrodes, which are constructed with a steel core, do indeed have an increased strength compared with copper or brass electrodes, however, all of these erosion electrodes have the disadvantage that if they have comparatively high strengths, their electrical conductivity is very low and merely reaches 8 Sxc2x7m/mm2. This comparatively high tensile strength is particularly advantageous for the cutting of high or thick workpieces or, however, also for the cutting of very small parts since high wire tensions are here demanded. However, these known wire electrodes have the disadvantage that their erosion performance is relatively low. In particular, in the case of cutting very small parts with erosion-wire diameters of up to 10 xcexcm, high wire strengths are demanded in order to reduce deflection and vibration of the wire to a minimum. Tungsten or molybdenum wires have been used up to now for the cutting of very small parts with wire thicknesses of 100 xcexcm and less, however, tungsten or molybdenum wires are extremely expensive to manufacture. Erosion wires with a steel core and a brass outer layer have been unable to be successful up to now for this purpose since they always, in comparison to tungsten wires, showed a lesser strength and a poorer erosion behavior.\nThe basic purpose of the invention is to provide high-strength erosion electrodes with a core of steel of the abovementioned type and a method for their manufacture, which electrodes have strengths which are 1800 N/mm2 and higher and, in addition, have a comparatively high conductivity of 10 Sxc2x7m/mm2 and more."}
-{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor memory capable of executing a read test for stored contents comparatively easily, and more particularly to a nonvolatile semiconductor memory such as an electrically writable/erasable flash memory.\n2. Description of the Background Art\nIn general, the production cost of a memory implies the total of the cost required for a wafer process, an assembly and a test. The cost of the test depends on how many chips can be tested per unit time by means of one tester. In order to reduce the test cost to produce a more inexpensive memory, accordingly, it is required that the prolongation of a test time caused by an increase in a storage capacity of the memory should be minimized even if the storage capacity of the memory is increased twofold to fourfold with an increase in the capacity.\nIn order to shorten a test time of the memory, the following should be implemented:\n(1) a reduction in an operation time required for write, read or the like;\n(2) development of a test pattern having a higher defect detecting capability; and\n(3) development of a test mode capable of carrying out write/read at a higher speed.\nReferring to (1), an increase in a speed of production has been required. Therefore, the test time tends to be shortened comparatively easily without a special contrivance through an enhancement in transistor performance by microfabrication and a reduction in a load capacity.\nReferring to (2), there are various test patterns having high defect detecting capabilities. Typically, a checker board pattern can be taken as an example.\nFIG. 95 is a diagram illustrating an example of the checker board pattern. FIG. 95 shows a checker board pattern having one bit/cell (binary-value).\nAs shown in FIG. 95, a checker board pattern CHK2 is a test pattern having a repetitive cycle of 2 bits in which adjacent bits always have a relationship of xe2x80x9c0xe2x80x9d and xe2x80x9c1xe2x80x9d. The checker board pattern CHK2 can detect open (disconnection) and a short circuit of a word line, open and a short circuit of a bit line, and a defective short circuit of floating gates in a nonvolatile semiconductor memory represented by a flash memory.\nFIG. 96 is a diagram illustrating a checker board pattern CHK4 having 2 bits/cell (quaternary-value). As shown in FIG. 96, 2-bit patterns CHK4-A to CHK4-D are repeated every 4-bit repetitive cycle in order to correspond to the quaternary-value (2 bits/cell) to be multiple-valued storage through application of the checker board pattern CHK2 shown in FIG. 95.\nFIG. 97 is a diagram illustrating a checker board pattern CHK8 having 3 bits/cell (octal-value). As shown in FIG. 97, a 3-bit pattern is repeated every 8-bit repetitive cycle in order to correspond to the octal-value (3 bits/cell) to be multiple-valued storage.\nIn the checker board pattern CHK4 and the checker board pattern CHK8, it is possible to detect a defective mode of the checker board pattern CHK2, and furthermore, to detect that all multiple-valued data can be written and read or not in the same word line and the same bit line.\nReferring to (3), there have been various methods. In general, if a test mode is incorporated into a chip, a chip area tends to be increased due to an incorporated circuit. Consequently, a cost required for a wafer process is increased. Accordingly, in the case in which the test mode is to be incorporated into the chip, it is necessary to note that the total production cost should be minimized.\nNext, the trend of a product of a flash memory and that of development will be described.\nAs an alternative to an EPROM (electrically writable nonvolatile semiconductor memory), a flash memory has spread for code storage. In recent years, a flash memory for mass data storage has spread more increasingly than the flash memory for code storage. In the case in which the flash memory for data storage carries out writing and reading randomly, it operates at a lower speed than that of the flash memory for code storage. However, in the case in which the flash memory for data storage carries out writing and reading sequentially, it can operate at a higher speed than that of the flash memory for code storage. The flash memory for data storage has had a larger capacity to exceed a DRAM through microfabrication of a processing pattern and a multiple-valued technique for storing multibit data in one memory cell.\nNext, a test time for the flash memory will be described.\nThe flash memory for code storage has such a structure that a reading operation can be carried out at a much higher speed than a writing operation. Therefore, a time required for a read test can be almost ignored with respect to the whole test time. However, the flash memory for data storage carries out the writing operation at a higher speed and has a larger capacity than the flash memory for code storage. Therefore, the time required for the read test cannot be disregarded with respect to the whole test time.\nFor example, a 256 Mbit flash memory has a 16 Ksector structure in which a 2 Kbyte write/read unit (hereinafter referred to as a xe2x80x9csectorxe2x80x9d) is present for 16K. Approximately 50 xcexcs is required for a reading head every sector (hereinafter referred to as a xe2x80x9c1st accessxe2x80x9d) and 50 ns is required for subsequent data transfer every byte.\nAccordingly, approximately 2.5 s ((50 xcexcs+50 nsxc3x972 Kbyte)xc2x716 Ksector) is required for carrying out a read test in the whole area of the 256 Mbit flash memory.\nIn a probing check stage of a wafer state, furthermore, it is hard to carry out a read test with a 2nd access 50 ns of a product specification due to a resistance and a capacitance of a probe needle and a resistance and a capacitance of a probe card, and a test time is further increased.\nIn general, the flash memory has an automatic writing/erasing function. The automatic writing function implies a function of repeating a write pulse applying operation and an operation (hereinafter referred to as a xe2x80x9cverify operationxe2x80x9d) for deciding whether desirable data are written (or erased) to (or from) an object memory cell for writing in accordance with a logic circuit (hereinafter referred to as a xe2x80x9ccontrol circuitxe2x80x9d) provided in an EEPROM, ending the repetition of the write pulse applying operation and the verify operation when it is decided that all the object memory cells store the desirable data and outputting a signal for giving, to the outside of the EEPROM, a notice that the writing operation (or the erasing operation) has been completed.\nIn order to decide whether xe2x80x9cthe contents stored in all the object memory cells are the desirable dataxe2x80x9d, an all latch deciding circuit (hereinafter referred to as an xe2x80x9cALL deciding circuitxe2x80x9d) is provided. The ALL deciding circuit serves to decide that all sense latches in a sense latch group provided for storing the result of read of the memory cell are xe2x80x9c1xe2x80x9d or xe2x80x9c0xe2x80x9d.\nFIG. 98 is a block diagram schematically showing a conventional ALL deciding circuit and a periphery thereof. In this specification, it is assumed that the case in which xe2x80x9c1xe2x80x9d is written to the memory cell is set to xe2x80x9cwritexe2x80x9d and the case in which xe2x80x9c0xe2x80x9d is written to the memory cell is set to xe2x80x9cerasexe2x80x9d. Referring to FIG. 98, description will be given in which a left memory cell group 31 is referred to as an L mat 31 and a right memory cell group 32 is referred to as an R mat 32.\nAs shown in FIG. 98, a sense latch group 33 is provided between the L mat 31 and the R mat 32. The sense latch group 33 transmits and receives data in a sector unit to and from the L mat 31 or the R mat 32. The latch data of the sense latch group 33 are output to an ALL deciding circuit 34.\nThe ALL deciding circuit 34 receives control signals LorR, 0or1 and ENABLE from an external control CPU 35 and outputs a decision result ALL34 to the control CPU 35. xe2x80x9c0xe2x80x9d/xe2x80x9c1xe2x80x9d of the LorR designates reading from the L mat 31/R mat 32, xe2x80x9c0xe2x80x9d/xe2x80x9c1xe2x80x9d of the 0or1 designates write verify/erase verify, and xe2x80x9c0xe2x80x9d/xe2x80x9c1xe2x80x9d of the ENABLE designates inactivity/activity of the ALL deciding circuit 34.\nFIGS. 99 to 102 are diagrams illustrating an operation principle of a conventional ALL deciding circuit, in which the ALL deciding circuit has a function of deciding that all sense latches are xe2x80x9c0xe2x80x9d.\nAs shown in FIGS. 99 to 102, the sense latch group 33 is provided between the L mat 31 and the R mat 32, and nodes N11 and N12 of each latch L33 of the sense latch group 33 are connected to the L mat 31 and the R mat 32 in one memory cell unit, respectively. The node N11 and the node N12 in the latch L33 are constituted to have a logical inverting relationship. A data latch group 36 is provided on the opposite side of the sense latch group 33 with the L mat 31 interposed therebetween, and a data latch group 37 is provided on the opposite side of the sense latch group 33 with the R mat 32 interposed therebetween.\nFirst of all, the ENABLE is set to xe2x80x9c1xe2x80x9d for initialization and the ALL deciding circuit 34 is thus brought into an active state.\nAs shown in FIG. 99, when the write verify for the L mat 31 is to be carried out, LorR=xe2x80x9c0xe2x80x9d (L) and 0or1=xe2x80x9c1xe2x80x9d are set so that data read from the L mat 31 are latched onto each latch L33 of the sense latch group 33. If xe2x80x9c1xe2x80x9d is normally written to the L mat 31, the nodes N11 of all the latches L33 of the sense latch group 33 are set to xe2x80x9c1xe2x80x9d and the nodes N12 are set to xe2x80x9c0xe2x80x9d.\nAccordingly, it is possible to execute the L mat write verify by deciding the state of the nodes N12 of all the latches L33 in the sense latch group 33 (whether all of them are xe2x80x9c0xe2x80x9d) through the ALL deciding circuit (R side sense latch decision).\nAs shown in FIG. 100, similarly, when the erase verify for the L mat 31 is to be carried out, LorR=xe2x80x9c0xe2x80x9d and 0or1=xe2x80x9c0xe2x80x9d are set. If xe2x80x9c0xe2x80x9d is normally written to the L mat 31, the nodes N11 of all the latches L33 in the sense latch group 33 are set to xe2x80x9c0xe2x80x9d and the nodes N12 are set to xe2x80x9c1xe2x80x9d.\nAccordingly, it is possible to execute the L mat erase verify by deciding the state of the nodes N11 of all the latches L33 in the sense latch group 33 through the ALL deciding circuit (L side sense latch decision).\nAs shown in FIG. 101, when the write verify for the R mat 32 is to be carried out, LorR=xe2x80x9c1xe2x80x9d (R) and 0or1=xe2x80x9c1xe2x80x9d are set so that data read from the R mat 32 are latched onto each latch L33 of the sense latch group 33. If xe2x80x9c1xe2x80x9d is normally written to the R mat 32, the nodes N12 of all the latches L33 of the sense latch group 33 are set to xe2x80x9c1xe2x80x9d and the nodes N11 are set to xe2x80x9c0xe2x80x9d.\nAccordingly, it is possible to execute the R mat write verify by deciding the state of the nodes N11 of all the latches L33 in the sense latch group 33 through the ALL deciding circuit (L side sense latch decision).\nAs shown in FIG. 102, similarly, when the erase verify for the R mat 32 is to be carried out, LorR=xe2x80x9c1xe2x80x9d and 0or1=xe2x80x9c0xe2x80x9d are set. If xe2x80x9c0xe2x80x9d is normally written to the R mat 32, the nodes N12 of all the latches 33 in the sense latch group 33 are set to xe2x80x9c0xe2x80x9d and the nodes N11 are set to xe2x80x9c1xe2x80x9d.\nAccordingly, it is possible to execute the R mat erase verify by deciding the state of the nodes N12 of all the latches L33 in the sense latch group 33 through the ALL deciding circuit (R side sense latch decision).\nThe ALL deciding circuit carries out the R or L side sense latch decision based on a logical expression of {(LorR) X OR (0or1)}. It is sufficient that the R side sense latch decision is carried out with xe2x80x9c1xe2x80x9d and the L side sense latch decision is carried out with xe2x80x9c0xe2x80x9d.\nBased on the result of the decision of the ALL deciding circuit for deciding whether all the memory cells in the sector are xe2x80x9c0xe2x80x9d, thus, it is possible to execute the write verify and erase verify operations without outputting read data for each bit from data input/output pins.\nIn this case, a time of approximately 150 xcexcs (50 xcexcs+50 nsxc2x72 Kbyte) required for normal reading per sector can be shortened to 50 xcexcs+xcex1 (xcex1 less than 1 xcexcs). Therefore, the read test time can be shortened to approximately one-third.\nIn this case, however, it is premised that all the write data in the sectors are identical. For this reason, there has been a problem in that the write data cannot be used for test reading through the checker board pattern CHK2, CHK4, CHK8 or the like which has a high defect detecting capability.\nA first aspect of the present invention is directed to a semiconductor memory comprising a plurality of memory cells, each of which can store N-valued (Nxe2x89xa72) information, a data reading device for reading a predetermined number of read data from a predetermined number of memory cells out of said plurality of memory cells during a reading operation for a test, and a deciding device for classifying the predetermined number of read data into K (Kxe2x89xa72) groups and outputting a decision result based on whether all the read data in the respective K groups are identical during the reading operation for a test.\nA second aspect of the present invention is directed to the semiconductor memory according to the first aspect of the present invention, wherein the K includes N, the memory cells include memory cells arranged in a matrix defined by first and second directions, the predetermined number of memory cells include memory cells provided in the same position in the second direction and provided in series in the first direction, and the deciding device classifies the predetermined number of read data such that the predetermined number of memory cells are classified into the same groups at N intervals in the second direction.\nA third aspect of the present invention is directed to the semiconductor memory according to the second aspect of the present invention, wherein the N-value includes a 2m-value (mxe2x89xa71).\nA fourth aspect of the present invention is directed to the semiconductor memory according to any of the first to third aspects of the present invention, the deciding device includes a sense storing device for sensing and storing the predetermined number of read data, and a decision result output device for deciding whether all the read data in the respective K groups are identical based on stored contents of the sense storing device and for outputting a result of the decision.\nA fifth aspect of the present invention is directed to the semiconductor memory according to any of the first to fourth aspects of the present invention, wherein the N-value includes a multiple-value to be ternary or more, the reading operation for a test includes first to Lth (Lxe2x89xa72) partial reading operations for a test which have different reading conditions and the read data include first to Lth 1-bit read data, the data reading device reads the predetermined number of first to Lth 1-bit read data during execution of the first to Lth partial reading operations for a test, the result of decision includes first to Lth partial decision results, and the deciding device outputs an ith (i=1 to L) partial decision result based on whether all ith 1-bit read data in the respective K groups are identical during an ith partial reading operation for a test.\nA sixth aspect of the present invention is directed to a semiconductor memory comprising a plurality of memory cells, each of which can store N-valued (Nxe2x89xa72) information, a data reading device for reading a predetermined number of read data from a predetermined number of memory cells out of said plurality of memory cells during a reading operation for a test, an expectation storing device for storing a predetermined number of expectation data, and a deciding device for outputting a result of decision based on a result of comparison of the predetermined number of read data with the predetermined number of expectation data during the reading operation for a test.\nA seventh aspect of the present invention is directed to the semiconductor memory according to the sixth aspect of the present invention, the deciding device includes a sense storing device for sensing and storing the predetermined number of read data, and a decision result output device for outputting the result of decision based on a result of comparison of stored contents of the sense storing device with stored contents of the expectation storing device.\nAn eighth aspect of the present invention is directed to the semiconductor memory according to the seventh aspect of the present invention, wherein the N-value includes a multiple-value to be ternary or more, the reading operation for a test includes first to Lth (Lxe2x89xa72) partial reading operations for a test which have different reading conditions, the read data include first to Lth 1-bit read data and the expectation data include first to Lth 1-bit expectation data, the data reading device reads the predetermined number of first to Lth 1-bit read data every execution of the first to Lth partial reading operations for a test respectively, the result of decision includes first to Lth partial decision results, the deciding device outputs an ith (i=1 to L) partial decision result based on a result of comparison of the predetermined number of ith 1-bit read data with the predetermined number of ith 1-bit expectation data during the ith partial reading operation for a test, and the predetermined number of second to Lth 1-bit expectation data are obtained by changing the predetermined number of first to (Lxe2x88x921)th 1-bit expectation data based on the stored contents of the expectation storing device and the sense storing device, respectively.\nA ninth aspect of the present invention is directed to the semiconductor memory according to any of the sixth to eighth aspects of the present invention, wherein the expectation storing device includes a data storing device for temporarily storing data when transmitting and receiving data between the memory cells and an outside.\nAs described above, according to the first aspect of the present invention, a comparatively complicated test pattern having the same value set to the K groups is written to the memory cells and the reading operation for a test is then executed to obtain the result of decision. Consequently, the read test for the memory cells can be carried out at a high speed.\nMoreover, it is sufficient that the deciding device has the function of deciding whether all the read data in the respective K groups are identical. Therefore, the circuit area of the semiconductor memory is rarely increased due to the provision of the deciding device therein.\nAccording to the second aspect of the present invention, a predetermined number of memory cells are classified into the same group at the N intervals in the second direction. Therefore, it is possible to carry out the read test in which the checker board pattern having the N-value is set to be a test pattern.\nAccording to the third aspect of the present invention, it is possible to carry out a read test in which a checker board pattern having a repetitive cycle of m bits is set to be the test pattern.\nAccording to the fourth aspect of the present invention, a predetermined number of read data are sensed and stored in the sense storing device. Consequently, it is possible to obtain a result of decision with high precision.\nAccording to the fifth aspect of the present invention, the result of decision comprising the first to Lth partial decision results is obtained. Consequently, it is possible to carry out the read test in a multiple-valued storage state for the memory cell without hindrance.\nAccording to the sixth aspect of the present invention, a predetermined number of expectation data are stored from the outside into the expectation storing device. Consequently, it is possible to carry out a read test based on an optional test pattern. Moreover, in the case in which the predetermined number of expectation data are utilized in common between the predetermined number of data read plural times, it is preferable that the predetermined number of expectation data should be stored in the expectation storing device at a first time. Therefore, a time required for storing the predetermined number of expectation data in the expectation storing device can be omitted during the reading operation of the predetermined number of read data at and after a second time. Correspondingly, the read test can be carried out at a high speed.\nMoreover, it is preferable that the deciding device should have the function of outputting the result of decision based on the result of comparison of the predetermined number of read data with the predetermined number of expectation data. Therefore, the circuit area of the semiconductor memory is rarely increased due to the deciding device constituted therein.\nAccording to the seventh aspect of the present invention, the predetermined number of read data are sensed and stored in the sense storing device. Consequently, it is possible to obtain a result of decision with high precision.\nAccording to the eighth aspect of the present invention, the result of decision comprising the first to Lth partial decision results is obtained. Consequently, it is possible to carry out the read test in a multiple-valued storage state for the memory cell without hindrance. In this case, the predetermined number of second to Lth 1-bit expectation data are obtained by changing the predetermined number of first to (Lxe2x88x921)th 1-bit expectation data based on the contents stored in the expectation storing device and the sense storing device, respectively. Therefore, it is sufficient that only the predetermined number of first 1-bit expectation data should be stored in the expectation storing device.\nAccording to the ninth aspect of the present invention, the expectation storing device does not need to be added specially for expectation data storage.\nAn object of the present invention is to provide a semiconductor memory capable of executing a read test at a high speed based on a comparatively complicated test pattern without increasing a circuit area. These and other objects, features, aspects and advantages of the present invention will become more apparent from the following detailed description of the present invention when taken in conjunction with the accompanying drawings."}
-{"text": "The invention relates to a paging receiver comprising a receiving circuit for the receiving and detecting of a first data block transmitted by a transmitting device, of which a second data block represents an address code for one or more receivers;\na first memory in which a received first data block is stored;\na first comparator which detects a first equality if the received address code is equal to an address code assigned to the receiver;\na second memory with two or more memory locations which are each suitable for the storing therein of a third data block of various first data blocks;\na first control circuit which, when the first equality occurs, scans locations of the second memory, compares by means of a second comparator the contents of corresponding sections of third data blocks which are present in the first memory and in the scanned location respectively, and alters the content of the scanned location depending on a second equality detected by the second comparator;\na second control circuit for the selection of a location of the second memory;\nand an indicating element for the indication of a section of the third data block which is present in a location selected by the second control circuit.\nA paging receiver of this type is known from British Pat. No. GB.2,101,779. In the known paging receiver the first, second and third data blocks are equal to each other and equal to the section of the third data block which is read out of the second memory and which can be displayed as a message on a window of the indicating element. On receipt of the data block the first control circuit will scan all the locations of the second memory and compare the content of each scanned location with the information content of the received data block. If an equality is detected in this process, the storage of the received data block in the second memory is prevented. If, after scanning all the locations, no equality is detected the oldest message present in the second memory is erased and the memory space which becomes available as a result of this is used for the storage of the newly received data block.\nIn the known paging receiver a newly received data block cannot be stored in a random location of the second memory. This is in particular troublesome if the oldest message recorded in the second memory contains more important information for the carrier of the paging receiver than the other messages recorded in the first and second memories. The importance of this drawback becomes particularly obvious in paging receiver systems with which emergency calls can be sent to service personnel such as in a nurse paging system.\nAnother important drawback of the known paging receiver is that a message no longer considered of importance by a user of the transmitting device is maintained in the second memory of the receiver so that the carrier of the receiver is inaccurately informed and may take an incorrect action.\nThe object of the invention is to eliminate the drawbacks of the known paging receiver."}
-{"text": "This invention relates to a device to be used for tending straps and also to be used as a whistle. More particularly, this invention relates to a device for tending straps with a whistle integrated within the device where both the whistling and strap-tending functions are independently operable.\nDevices for tending straps take many forms. The buckle on a typical belt is maybe the most recognizable, but fasteners with similar functions encompass a wide range of designs and related functions, see U.S. Pat. No. 5,794,316 for example. The xe2x80x9ctending strapxe2x80x9d function is intended to encompass the various functions performed by devices that are attached to straps, or more specifically to belts, webbing, or cords, and that work to secure the straps under tension (a buckle or webbing adjuster for example), confine them to a particular area (a slider), or limit movement of other items along the strap (a cord lock). The articles utilizing these designs are ubiquitous where one finds hikers, boaters, campers, climbers, and other persons employing similar gear. Many of these activities give rise to emergencies or other situations where the spoken word either does not carry far enough, is not loud enough to be heard over surrounding noise, or is not distinct enough to draw attention. It is often recommended that the participant carry a whistle for safety in activities where emergency situations can exist. Such needs for help, however, tend to be rather unexpected. This unexpected quality results in people not equipping themselves with whistles at all times and thus not having one when such an emergency occurs. If the whistle could be incorporated into an otherwise useful item then a person would be more likely to have one available at all times, particularly during emergencies. Aside from emergency uses, a whistle may be merely a handy means for the user to notify others of his or her presence, or otherwise get attention. There is thus a need for a device that has a function with a utility related to the activity that also incorporates a whistle that can be used relatively quickly when the need arises.\nWhistles in general take many forms, for examples see U.S. Pat. Nos. 5,546,887 and 5,507,246. Such whistles undeniably perform the necessary function if they are available. But such whistles are often not available for any number of reasons. A solitary whistle is not an item that most would habitually carry. Therefore a potential user would probably need a reason to include a whistle in the gear for that day. Most of the obvious reasons are related to emergencies and most people do not anticipate having emergencies. This would make a whistle a low priority item on anyone\"\"s packing list. And because emergencies are fortunately rare, the whistle would in fact generally be just another item cluttering up a purse, key-chain, or related devices. This limited need makes it easy and understandable to simply forget to carry a whistle even should one think it a good idea in general.\nBut many other things are carried constantly and without additional effort. Many of the activities giving rise to the need for a whistle have a particular type of gear associated with that activity. Although the various activities do not necessarily have gear in common, the individual items of gear have common elements. One group of common elements comprises the fasteners used to secure webbing, belts, straps, and cords on this gear. These devices range from your typical belt buckles to the more high tech molded side release buckles, tension adjusters, sliders, and cord locks. They are present in gear ranging from book bags and backpacks, to key-chains, bike helmets, and the cords on the hoods of many jackets.\nBuckles have been proposed incorporating tools such as a whistle, however the design has been unsatisfactory and has not gained wide spread acceptance. For example, U.S. Pat. Nos. 3,885,250 and 3,903,547 disclose a buckle with a prong that may be fashioned into a whistle. While specific details of how such a whistle might be formed are not disclosed, the general arrangement is not satisfactory because, being on a prong, such a whistle would be hidden and even if known, use would require disengagement of the retained strap or belt. Additionally, the strap must be large enough to house the whistle. This requirement prevents incorporation of a whistle where the strap is of small cross-section because the prong must be accommodated within either a thick strap or two layers of strap. This also means that the strap must be specially manufactured to accept the flange, increasing manufacturing costs and reducing the application of the buckle because it must necessarily be part of a matched set to function. Finally, since the prong is rigid and extends into the strap the areas where the buckle can be located are limited to relatively flat locations that extend the length of the prong and buckle.\nThus there exists a present need for a device with an integral whistle that can function with a variety of strap sizes and shapes to fasten, tension, lock or generally tend the straps and where the whistle is unobtrusive yet apparent and accessible to the user.\nThe preferred embodiments of the present invention combine the function of a whistle with that of a device for tending straps, cords, belts, or webbing of various shapes and sizes. The whistle is combined with the strap-tending device in a manner that allows the whistle to be used even when the device is tending the straps as designed. Given that a strap runs in a general direction and that the strap-tending device typically performs a function that is oriented along the same general direction, the placement of the whistle is preferably transverse to that direction to facilitate access to the whistle, decrease device dimensions, and improve its manufacturability. The strap-tending functions of the preferred embodiments include: a two-piece side release buckle with the whistle incorporated in the female half and the strap adjusting mechanism in the male half; a single piece strap adjuster with the whistle located opposite the anchor bars and slots from the extension designed to facilitate releasing tension; a cord lock with the whistle integrated into the barrel; and a slider with the whistle integrated to one side of the slots designed to engage the straps. In all cases the whistle is accessible and functional with the intended strap or straps in place and the device performing its intended strap-tending function."}
-{"text": "Advances in semiconductor manufacturing technology have resulted in, among other things, reducing the cost of sophisticated electronic products to the extent that the integrated circuits fabricated in accordance with these semiconductor manufacturing technologies have become ubiquitous in the modern environment.\nAs is well-known, integrated circuits are typically manufactured in batches, and these batches usually contain a plurality of semiconductor wafers within and upon which integrated circuits are formed through a variety of semiconductor manufacturing steps, including, for example, depositing, masking, patterning, implanting, etching, planarizing and so on.\nCompleted wafers are tested to determine which die, or integrated circuits, on the wafer are capable of operating according to predetermined specifications. In this way, integrated circuits that cannot perform as desired are not packaged, or otherwise incorporated into finished products.\nIt is common to manufacture integrated circuits on roughly circular semiconductor substrates, or wafers. Further, it is common to form such integrated circuits so that conductive regions disposed on, or close to, the uppermost layers of the integrated circuits are available to act as terminals for connection to various electrical elements disposed in, or on, the lower layers of those integrated circuits. Such conductive regions are commonly referred to as pads. Pads are commonly used to provide electrical access to the integrated circuit both during testing and during the operation of the integrated circuit as it is incorporated into a final product. Typical pads are formed from aluminum. It is well-known that the surfaces of aluminum which are exposed to the atmosphere will oxidize. These oxidation layers interfere with the formation of low resistance electrical connection to the pads.\nDuring the process of testing the performance of the integrated circuits, the pads are commonly contacted with probe needles, or other contact structures, of a probe card. It is through such temporary connections that a test apparatus may electrically interact with an integrated circuit.\nConventionally, the probe needles, or other contact structures, of the probe card are brought into physical contact with the pads and then moved laterally to \u201cscrub\u201d the pad. The scrubbing operation is intended to break through the oxide layer on the top surface of the pad, thereby providing reduced contact resistance. Unfortunately, scrubbing disturbs the pad structure and can contribute to yield loss due to failure of bond wires to properly attach to the disturbed pad structure.\nWhat is needed are methods and apparatus for providing low contact resistance connections to pads of integrated circuits without disturbing or substantially redistributing the material from which the pads are formed."}
-{"text": "FIG. 1 is a cross-section of a known formation testing/sampling probe assembly 100 viewed in a plane containing the wellbore axis. The assembly 100 includes a packer 101 configured to be pressed against a formation of interest via hydraulically actuated pistons schematically shown as 103 and 104. Deformation of the packer 101 around an extendable probe barrel 102 and against the formation creates a seal that isolates the formation fluids and pressure from the wellbore environment. The probe barrel 102 extends forward, due to high pressure hydraulic oil entering chamber 105 at the same time pistons 103 and 104 are actuated to apply an extension force. The area of investigation by the probe assembly 100 is limited to the area that is in direct contact with the extendable probe barrel 102, usually called the orifice. This area is limited by the diameter of the probe barrel 102, which is about the same dimension as the diameter of the hole 106 in the packer 101 through which the probe barrel 102 passes.\nOther prior art has sought to improve upon the design of the probe assembly 100 by increasing the diameter of the probe barrel. Nonetheless, the area of investigation is still limited to the largest diameter of the probe barrel, which is currently between two and three inches. This could be problematic for formations with thin laminations where the formation may have a small (e.g., \u02dc0.5 inches thick) production zone sandwiched between thicker zones of impermeable formation. With prior art designs, finding a non-producing zone is much more likely due to the small lateral extent (the probe barrel diameter) of the area of investigation.\nOne attempt at addressing this thin lamination problem involved elongating the packer, such as with the known probe assembly 200 shown in FIG. 2. The elongated packer 201 of the probe assembly 200 is pressed against the formation by backing plate 202 to create a seal. However, in this case, an annular shaped metal spacer 207a fills the upper hydraulic chamber 207, preventing the probe barrel 208 from extending forward. Since the packer 201 and backing plate 202 are elongated along the central axis of the wellbore, a number of formation laminations could be investigated. The axial length of investigation using the probe assembly 200 could approach several inches, including one known embodiment in which the length was about seven inches.\nHowever, problems with the probe assembly 200 have been discovered. For example, the packer 201 is not as constrained at its inner boundary as it is near the raised rim 203 of the backing plate 202 that defines the orifice of the probe. It has been found that the packer material 201 migrates into the probe's orifice during formation testing operations. The large difference in pressures between the wellbore, pushing on surfaces 204, and the probe orifice, pushing on the packer surface near the raised rim 203, forces the packer material to move into the orifice. This movement decreases the ability of the packer 201 to seal against the formation and prevents the tool's pressure gauges from equilibrating to that of the pressure of the formation fluid. As the elastomeric material of the packer 201 is sucked into the probe orifice, the hydraulic pressure declines in the pistons (schematically shown in FIG. 2 as 205 and 206) applying the extension force (a consequence of the tool design). This, in turn, results in a decline in the force pushing the packer against the formation, further exacerbating the situation, and resulting in both packer damage and loss of seal. This design has lead to a high number of lost seals.\nAnother problem that exists with the design shown in FIG. 2 is that the upper surface of the raised rim 203 of the backing plate 202 is located well below the surface of the elastomer 201 that seals against the wellbore wall. As a consequence, elastomer must be compressed before the raised rim 203 makes contact with the wellbore wall. This is of considerable disadvantage because it requires performing large volume pretests. On the other hand, if this difference in heights between the surface of the raised rim 203 and the sealing surface of the elastomer 201 didn't exist, then it would not be possible for the elastomer 201 to form a seal against the wellbore wall, because the metal spacer 207a makes the probe barrel 208 and the backing plate 202 move in unison."}
-{"text": "In subsea offshore well operations, a subsea installation may include hydraulically and also electrically controlled equipment, such as valves which are remotely operated from a platform structure located at the surface of the ocean. Hydraulic and electrical control lines extending from the platform to the subsea installation are desired to be releasably connected to the subsea installation so that if the platform must move to another location or be released from its connection to the subsea station, such release may be readily made. In some instances, such a control line has been associated with tension line means integral therewith so that the line, sometimes called an umbilical line, may be locked at its lower end to the subsea installation.\nAn example of such an umbilical line construction is shown in my copending application Ser. No. 771,799 owned by a common assignee wherein the umbilical line extends from the subsea installation to a buoyant structure located about 300 feet below the sea surface and the umbilical line continues above the buoyant structure to a floating platform. The buoyant structure is connected with a marine riser system having a lower riser portion and an upper riser portion so that the upper riser portion and the upper portion of the umbilical line can be disconnected at the buoyant structure and the floating platform moved to another location upon such disengagement. The buoyant structure and the lower riser portion are free-standing to permit reengagement of the upper riser portion at a later time. It is desirable that the buoyant structure be securely tethered to the subsea installation so that in the event the lower riser portion should break, the buoyant structure, which exerts an upward tension on the lower riser portion, will not be suddenly released and rapidly travel to the surface and perhaps strike and cause damage to the floating platform normally thereabove.\nIn prior proposed locking systems for guide line connectors for use in such installations, manually operated locking devices were maintained in locked position by mechanical springs and were unlocked by use of a submarine or by use of go devil means."}
-{"text": "Earrings are not easily stored in any organized fashion because of their small size and odd shapes, and the fact that they are used in pairs. The usual storage is accomplished by simply keeping them loose in a box or drawer, with the result that it is difficult to find a matching pair.\nDisplay racks found in retails stores usually involve displaying small cards or boxes in which are fastened a single pair of earrings in a manner such that the earrings are not readily detachable. As such, these racks are totally unsuitable for use in the home where one wants to quickly find and remove from storage a selected pair of earrings. Insofar as is known, there is no such storage rack available at the present time, nor has one been disclosed in the prior art.\nIt is an object of this invention to provide an improved holder for earrings for use in the home. It is another object to provide a holder for earrings in which the earrings are easily pinned to a fabric. Still other objects will become apparent from the more detailed description which follows."}
-{"text": "Interests in energy storage technologies have been increasingly higher recently. As applications are expanded to energy of mobile phones, camcorders and notebook PCs, and furthermore, to electric vehicles, efforts on the research and development of electrochemical devices have been more and more materialized.\nElectrochemical devices are fields receiving most attentions in such aspects and among these, development of secondary batteries capable of charge and discharge have been the focus of attention, and in developing such batteries, research and development on the design of new electrodes and batteries for enhancing capacity density and energy efficiency have been recently progressed.\nAmong currently used secondary batteries, lithium secondary batteries developed in early 1990s have received attentions with advantages of having high operating voltage and significantly higher energy density compared to conventional batteries such as Ni-MH, Ni\u2014Cd and sulfuric acid-lead batteries using an aqueous liquid electrolyte.\nA lithium secondary battery has a structure of an electrode assembly including a positive electrode, a negative electrode and a separator interposed between the positive electrode and the negative electrode being laminated or wound, and is formed by embedding the electrode assembly in a battery case, and injecting a non-aqueous liquid electrolyte thereinto. The lithium secondary battery produces electric energy through oxidation and reduction reactions occurring when lithium ions are intercalated/deintercalated in the positive electrode and the negative electrode.\nIn a common lithium secondary battery, a negative electrode uses lithium metal, carbon and the like as an active material, and a positive electrode uses lithium oxides, transition metal oxides, metal chalcogen compounds, conductive polymers and the like as an active material.\nAmong these, a lithium secondary battery using lithium metal as a negative electrode mostly attaches lithium foil on a copper current collector, or uses a lithium metal sheet itself as an electrode. Lithium metal has low potential and high capacity, and has received much attention as a high capacity negative electrode material.\nWhen using lithium metal as a negative electrode, electron density non-uniformization may occur on the lithium metal surface when operating a battery due to various reasons. As a result, a branch-shaped lithium dendrite is produced on the electrode surface causing formation and growth of projections on the electrode surface, which makes the electrode surface very rough. Such lithium dendrite causes, together with battery performance decline, separator damages and battery short circuits in severe cases. As a result, a temperature in the battery increases causing a risk of battery explosion and fire.\nIn order to resolve such problems, researches such as introducing a polymer protective layer or an inorganic solid protective layer to a lithium metal layer, increasing a concentration of a salt of a liquid electrolyte, or using proper additives have been urgently required."}
-{"text": "1. Field of the Invention\nThe present invention relates to the field of leakage coaxial cables which are used for radio communication in closed in areas such as the interior of buildings, tunnels or underground markets, where no ordinary radio waves can be received.\n2. Description of the Prior Art\nA leakage coaxial cable typically has slots formed in the exterior thereof which are spaced at predetermined intervals along the conductor, so that electromagnetic waves propagating inside the coaxial cable are partially radiated from the conductor into an external space through the slots.\nWhen a fixed signal source is connected to the leakage coaxial cable, the signal is radiated into the external space and may be received by a mobile station running near the leakage coaxial cable. In addition, a signal transmitted by the mobile station may be received by the fixed station through the leakage coaxial cable.\nA common application for leakage coaxial cables has been in disaster warning and prevention system for use in building, tunnels and underground markets. It is important in such applications to make certain that the cables are fireproof. The ability of prior art leakage coaxial cables to withstand high heat has been limited. Consequently, there has existed a long and unfilled need in the prior art for a method of making leakage coaxial cables that have transmission characteristics which will not degrade during a fire related emergency.\nA conventional leakage coaxial cable has an external conductor with slots formed therein to radiate an electromagnetic wave which is propagating inside the cable outwardly into an external space. The external conductor is disposed coaxially around an internal conductor with an insulating member therebetween. The external conductor is covered by a protective sheath. In order to minimize the transmission loss of the electromagnetic wave, the insulating member is preferably made of a low-loss plastic material, such as polyethylene or polystyrene. The external conductor is preferably made of high-conductivity material such as alminum or copper. A polyester film is laminated on the external conductor with adhesive to compensate for the decrease of the mechanical strength of the conductor which is caused by the formation of the slots. The protective sheath is preferably made of polyethylene or polyvinyl chloride.\nIf a leakage coaxial cable thus constructed encounters a fire, the protective sheath will burn away and the external conductor will be directly exposed to the flames. The polyester will burn, and the plastic insulating member will melt. The molten plastic will flow through the slots formed in the external conductor, ignite and drop from the cable while burning. The burning molten plastic may actually contribute to the spread of the fire and may burn the skins or clothes of persons fighting the fire or running away from the fire.\nOne example of a prior art method of making a leakage coaxial cable refractory has been disclosed in Japanese Utility Model Application Publication No. 16682/1977. In that method, a heat-resistant tape which is made of an inorganic material such as asbestos is spirally wound between a polyethylene insulating member and an external conductor. The heat-resistant tape prevents the melting of the polyethylene insulating material for a long time. In addition, the heat-resistant tape maintains the insulation between the internal and external conductors even after the polyethylene insulating member has been molten. Therefore, the radio communication properties of the cable can be maintained unchanged for some time after the occurrence of the fire.\nHowever, the leakage coaxial cable disclosed in Japanese Utility Model Application Publication No. 16682/1977 is still not an ideal solution to the problem discussed above, since the polyethylene insulating material will eventually ooze out through the stitch lines of the non-organic tape or the seams of the spirally wound tape and flow out through the slots. In addition, since the tape is relatively thick, on the order of 0.25 to 0.5 mm, the dielectric value between the internal and external conductors is large, causing large transmission losses.\nAnother example of a prior art conventional leakage coaxial cable is disclosed in Japanese Utility Model Application (OPI) No. 3537/1980. In that cable, a heat-resisting tape of polyimide resin is spirally wound between an external conductor and a polyethylene insulating member. Because of the presence of the heat-resisting tape, the internal and external conductors are not short-circuited even if the polyethylene insulating member is molten.\nHowever, that leakage coaxial cable also has the problem of the polyethylene insulating member oozing out through the seams of the spirally wound tape, flowing out through the slots and dropping from the cable while burning."}
-{"text": "Abbreviations used herein shall have the following meanings:\nAF Amplify-and-Forward\nCRC Cyclic Redundancy Check\nDCI Downlink Control Information\nDF Decode-and-Forward\nDFTS Decode and Forward TS\nDL Downlink\nFSR Frequency Selective Relay (Repeater)\nFTR Frequency Translating Relay (Repeater)\nOFDM Orthogonal Frequency Division Multiplex\nOFR On-Frequency Relay (Repeater)\nPDCCH Physical Downlink Control CHannel\nCRNTI Cell Radio Network Temporary Identifier\nTTI Transmission Time Interval\nUL Uplink\nOne recent development within modern telecommunication standards is the so-called Long Term Evolution (LTE) radio interface, and a further development of the LTE, namely LTE Advanced. These are both OFDM based systems.\nOne of the most important improvement areas in the so-called LTE-Advanced technology is the increase of data rates available for users at the cell edge and indoor. A very promising technique to achieve high data rates in such difficult locations is the deployment of relays. Relays are usually classified into Layer 1 (L1), Layer 2 (L2), and Layer 3 (L3) relays depending on which Open System Interconnection (OSI) layer they operate on. The OSI model is a conceptual model for telecommunication consisting of 7 layers (physical, data-link, network, transport, session, presentation, and application) Note however that the different layers refer to the user plane of the relay node and a L1 relay may use e.g. L3 control plane signaling.\nL1 relays are commonly denoted Amplify-and-Forward (AF) relays or sometimes equivalently repeaters. An AF repeater operates in the physical layer and its basic functionality is, as the name suggests, amplifying and then forwarding the received signal, including any received noise and interference.\nL2 relays operate in the data link layer and have the ability to detect and possibly correct errors that have occurred in the physical layer. L2 relays are therefore commonly called Decode-and-Forward (DF) relays as they decode the received data prior to retransmission. DF relays will, at the expense of an increased delay, not forward noise and interference.\nL3 relays operate in the network layer and are by the Third Generation Partnership Program (3GPP) regarded as being equivalent to eNodeBs (eNBs) that are wirelessly connected to a donor cell via self-backhauling. L3 relays have the same characteristics as L2 relays in the sense that they do not forward noise and interference as they perform decoding and error correction of the received signal prior to retransmission.\nThere are several known methods of utilizing various repeaters or relays to further improve the quality of performance in telecommunication systems. Some of the most common include:\nCooperative relaying which enables multiple relays to cooperate during transmissions to users. For example, the cooperation may be used for increased diversity or multiplexing of data.\nMulti-hop relaying which enables signals to be conveyed from a source to a destination over two or more wireless hops. The multiple hops are achieved by relaying the signal via one or more relay(s)/repeater(s). It may be used to reduce the end-to-end path loss and thus extending the coverage.\nOn frequency relays (or repeaters) (OFR) which are relays (repeaters) that forward on the same frequency band occupied by the received signal.\nFrequency translating relays (or repeaters) (FTR) which are relays (repeaters) that translate the retransmitted signal to another frequency band that is different from that occupied by the received signal.\nFrequency selective relays (or repeaters) (FSR) which are relays (repeaters) that may dynamically retransmit coordinated parts of the received signal bandwidth.\nThe increased path gain that comes from splitting of the signal path in two hops by repeating or relaying in an intermediate node brings several benefits: Data rates can be significantly increased; transmit power can be reduced and inter-cell interference falls rapidly. A multi-hop solution based on AF repeaters has some interesting characteristics compared to other DF relaying solutions. Since a repeater can receive and transmit on the same radio resource, which is not possible for DF types of relays, it is possible to operate without any duplex coordination loss between the two hops. A decode-and-forward (DF) relay can forward the data on the same frequency resource. However, since the decoding operation will result in an unavoidable delay the forwarding must take place at a later time instance, i.e. on another radio resource. In contrast, an amplify-and-forward (AF) repeater has a delay that typically is negligible compared to the transmission time interval, hence it can forward on the same radio resource. Repeaters also introduce less delay than DF relays which is beneficial for the performance of higher layer protocols such as TCP. Furthermore, a repeater is a simple device that typically is fairly cost efficient.\nIn particular, the use of OFR in OFDM based systems is interesting if the delay of the repeater is limited to the length of cyclic prefix of the OFDM modulation. In the air, the repeated signal path and the direct signal path add in the same way as normal multi-path does. In case of LTE, the additional time dispersion induced by the repeater does not result in any additional receiver complexity and/or reduced performance due to increased self interference as long as the total time dispersion is limited to the length of the cyclic prefix. Note that is not the case for single-carrier systems without cyclic prefix e.g. HSPA, where additional time dispersion typically increases the receiver complexity (i.e. more rake fingers are required) as well as the self-interference (i.e. signal components with a relative delay difference are non-orthogonal).\nDespite the benefits of utilizing AF repeaters, there are a few disadvantages that prevent the use from providing the full benefit of them.\nOne drawback with repeaters compared to DF relays is that they forward not only signals but also noise and interference. Furthermore, a major challenge for on-frequency repeaters is to sufficiently suppress the self-interference they induce.\nRepeaters (and relays) are efficient for both providing coverage in areas without coverage (see upper part of FIG. 1) and also to provide increased data rates to areas with weak signal strength (lower part of FIG. 1). This distinction is important since in the data-rate extension case the users will receive both a direct signal part as well as a repeated signal path, witch in the case of DF relaying will interfere with each other and in case of AF repetition and OFDM will add like multi-path. In addition, since AF repeaters amplify noise and interference they are only beneficial in case it is possible to replace one weak radio link with two significantly better radio links. This is more likely to be possible when the original radio link is weak due to some obstacle (e.g. a wall) that hinders the radio waves rather than pure propagation distance.\nIn addition, in the data rate extension scenario it is possible to dynamically turn the repeater on and off without losing coverage. It is also possible to do frequency selective repetition in the data-rate extension case without destroying the communication on the uplink and downlink control channels. Furthermore, the data-rate extension case is also particularly relevant for LTE and LTE-Advanced since, in order to compete with HSPA, their main business case is to provide high data-rates, which are only achievable in case the signal strength also is high.\nIn the coverage extension scenario, the options when it comes to advanced repeater behavior are more limited. It is not possible to e.g. turn the repeater off even when the repeater does not serve any UEs since that would leave the area with no coverage. In that case a UE wanting to perform an initial access would not be able to read the broadcast channel (BCH) and the system information blocks (SIBs) required for random access. Furthermore, an idle UE would not detect any paging messages sent from the network. In the coverage area extension case it is also not possible to perform frequency selective repetition on the downlink since that would hinder the UE from receiving the physical downlink control channel (PDCCH) that covers the whole downlink bandwidth. Also any frequency selective operation by the repeater on the uplink band must assure that the resources used for the physical uplink control channel (PUCCH) as well as the physical random access channel (PRACH) are always repeated.\nConsequently, there is a need for a more efficient use of repeaters."}
-{"text": "1. Field of the Invention\nThe present invention relates to a fixing apparatus provided in an image forming apparatus, such as a copying machine, a printer, and a facsimile, and a roller used in the fixing apparatus, and in particular, a fixing apparatus of an image forming apparatus and a roller using an induction heating system.\n2. Description of the Background\nIn a fixing apparatus that adopts an induction heating system and is used in an image forming apparatus, such as an electrophotographic copying machine or printer, there is a device that increases the fixing speed by speeding up a warming up time of the fixing apparatus. For example, JP-A-2002-295452 discloses a heating device that heats a metal sleeve, which is positioned on the outer periphery of an elastic layer of a heat roller and has a small heat capacity, using an induction coil to thereby shorten a warming-up time.\nHowever, in the known device, handling of expansion or contraction that occurs to the elastic layer when warming up or cooling the elastic layer is not mentioned.\nTherefore, in a case when a metal belt having a metal layer on the outer periphery of an elastic layer is provided, it is preferable to develop a fixing apparatus of an image forming apparatus capable of increasing the life of the metal belt and an elastic roller by preventing the metal belt or the elastic layer from being broken at an early stage regardless of expansion or contraction occurring to the elastic layer."}
-{"text": "Infrared (IR) thermal cameras can be used in a number of different situations, for example, when inspecting or surveying complex electrical systems such as transformers, switchgears etc., or water carrying systems such as heat exchangers, radiators etc. IR cameras are used for capturing and storing thermal radiation data. This thermal radiation data may then be displayed and viewed as thermal images and analyzed in order to, for example, find faulty electrical wirings or couplings, leaking water pipes, etc.\nHowever, various procedures and methods are being used in order to properly analyse the thermal radiation data and/or the thermal images of the IR camera, and these are not necessarily particularly intuitive and easily understandable by a user of the IR camera. The analysis of the thermal radiation data and/or the thermal images of an IR camera may also be a time-consuming task and may thus preclude a user of an IR camera from making decisions, predictions, and/or recommendations to clients while being on site and performing the IR imaging."}
-{"text": "Experimental and clinical observations have supported the concept that the hypothalamus plays a key role in the regulation of adenohypophysial corticotropic cells' secretory functions. Over 25 years ago it was demonstrated that factors present in the hypothalamus would increase the rate of ACTH secretion by the pituitary gland when incubated in vitro or maintained in an organ culture. However, a physiologic corticotropin releasing factor (CRF) was not characterized until ovine CRF (oCRF) was characterized in 1981. As disclosed in U.S. Pat. No. 4,415,558, the disclosure of which is incorporated herein by reference, oCRF was found to be a 41-residue amidated peptide. oCRF lowers blood pressure in mammals when injected peripherally and stimulates the secretion of ACTH and .beta.-endorphin.\nRat CRF (rCRF) was later isolated, purified and characterized; it was found to be a homologous, amidated hentetracontapeptide as described in U.S. Pat. No. 4,489,163, the disclosure of which is incorporated herein by reference. The formula of human CRF has now been determined to be the same as that of rCRF, and the terms rCRF and hCRF are used interchangeably.\nA CRF analog was subsequently developed having a high alpha-helical forming potential which is also a 41-residue amidated peptide. It is commonly referred to as AHC (alpha-helical CRF) and is described in U.S. Pat. No. 4,594,329, the disclosure of which is incorporated herein by reference. Other analogs containing D-isomers were developed, such as those shown in U.S. Pat. No. 5,278,146.\nSynthetic rCRF, oCRF and AHC stimulate ACTH and .beta.-endorphin-like activities (.beta.-END-Li) in vitro and in vivo and substantially lower blood pressure when injected peripherally. Antagonists of these compounds are disclosed in U.S. Pat. No. 4,605,642, issued Aug. 12, 1986, the disclosure of which is incorporated herein by reference. Cyclic CRF analogs exhibiting biopotency have been developed as disclosed in U.S. Pat. No. 5,245,009 (Sep. 14, 1993) and in pending U.S. patent application Ser. No. 78,558, filed Jun. 16, 1993, now U.S. Pat. No. 5,493,006.\nSince the foregoing discoveries, the search for improved CRF analogs has continued."}
-{"text": "In the UMTS (Universal Mobile Telecommunications System) network, the specifications of long term evolution (LTE) have been drafted for the purpose of further increasing high speed data rates, providing lower delays and so on (see non-patent literature 1). In LTE, as multiple-access schemes, a scheme that is based on OFDMA (Orthogonal Frequency Division Multiple Access) is used in downlink channels (downlink), and a scheme that is based on SC-FDMA (Single Carrier Frequency Division Multiple Access) is used in uplink channels (uplink). Also, successor systems of LTE (also referred to as, for example, \u201cLTE-advanced\u201d or \u201cLTE enhancement\u201d (hereinafter referred to as \u201cLTE-A\u201d)) have been developed for the purpose of achieving further broadbandization and increased speed beyond LTE, and the specifications thereof have been drafted (Rel. 10/11).\nIn relationship to LTE-A systems, a HetNet (Heterogeneous Network), in which small cells (for example, pico cells, femto cells and so on), each having local a coverage area of a radius of approximately several tens of meters, are formed within a macro cell having a wide coverage area of a radius of approximately several kilometers, is under study. Also, in relationship to HetNets, a study is in progress to use carriers of different frequency bands between macro cells and small cells, in addition to carriers of the same frequency band.\nFurthermore, for future radio communication systems (Rel. 12 and later versions), a study is in progress to introduce a new mechanism for small cell discovery. To be more specific, a user terminal detects the detection/measurement signal that is transmitted from a small cell, and establishes synchronization, conducts measurements and so on. Note that this detection/measurement signal may be referred to as the \u201cdiscovery reference signal\u201d (DRS). When a channel state measurement signal (for example, the CSI-RS (Channel State Information Reference Signal) is included in the DRS, the user terminal detects and measures this channel state measurement signal based on the timing of synchronization."}
-{"text": "1. Field of the Invention\nThis invention relates to a peptide which has an affinity for gp120, HIV (human immunodeficiency virus) envelope protein.\n2. Description of the Related Art\nThe therapy for HIV infection is usually chemotherapy, such as the nucleotide derivative AZT (3\u2032-azido-3\u2032-deoxythmidine). This AZT therapy or protease inhibitor, which was later developed, prolongs the life of HIV patients, but there are some problems, these are derived from the chemotherapy itself.\nThe problems are shown as follows: The first is chronic poisoning due to long term administration, the second is the appearance of an HIV virus resistant to the medicine during the therapy, the third is the appearance of malignant tumors in prolongation of the HIV patient's life, the fourth is that the recovery of the immune system can not be monitored, the fifth is that there is not a method to monitor treatment effect, etc. Since such chemotherapy is not basic therapy for HIV infections, most people anticipate the development of a vaccine.\nGenerally, the vaccine is an inactive treatment (in active vaccine) of a microbe of viruses; a weak activity virus which loses pathogenesis or a pseudo virus (live vaccine) which has no fatal effects to humans. However, although HIV itself is natively a weak activity virus, it is well known to stay long after having once invaded the body. In addition, the host cell of HIV is mainly a lymphocyte, which controls the immune system; furthermore, HIV spreads over hemophilic patients through blood-preparation. From these finding, even if it is assumed that we selected either an in-active or a weak vaccine, the development of an HIV vaccine has many problems with safety.\nAccordingly, an HIV vaccine is being developed which utilizes a part of the HIV envelope protein and inhibits further infection.\nFrom such an idea, many researchers performed an epitope analysis of gp120 in the HIV enveloped protein, and then, watched the V3 region (3rd hypervariable region) of gp120 as an epitope. But it was a true hypervariable region [Palker T. J., et al., Proc. Natl. Acad. Sci. USA 85:2709-2713, 1988; Rusche J. R., et al., ibid 85:3198-3202, 1988; Gouddsmit J., et all., ibid 85:4478-4482.1988; Matsushita S., et al., J. Virol. 62:2107-2114, 1988]. After this, a vaccine which used a part of this region as antigen was administrated to an HIV infected monkey as an infection inhibitory experiment, but the effectiveness has not yet been reported.\nAs well as this, Tam et al. devised further antigenecity for the above-mentioned peptide antigen (Tam et. al., Japanese patent publication (Tokuhyo) No. H 3-503539), but have not yet had success because in most parts of the V region, particularly in the V3 which is a convenient region for antibody preparation, mutation or deletion occurs.\nIn addition, a neutralized antibody, which inhibits the infection against lymphocyte, is developed. For example, in Japanese Patent Application No. 63-171385, after the production of a monoclonal antibody by using a part of the above mentioned peptide as antigen, a method is reported, which produces anti HIV chimera antibodies on hybridization due to genetic engineering at the level of the protein as the Fab\u2032 itself. But, although with such neutralizing antibodies it is possible to inhibit HIV infection to the lymphocyte at laboratory level, an antibody that can be used practically has not yet been developed.\nAs mentioned above, chemotherapy has some problems; drug tolerance in the virus and side effects in the host, another idea to solve the problem of removing the virus from the body is by plasmapheresis. Although this method to remove the HIV virus by using a pore size membrane filter (smaller than virus size) for plasmapheresis has been definitely proposed it is not yet possible to make a uniform pore size membrane. It is also possible that the pores will become clogged during plasmapheresis resulting in the deterioration of the membrane due to pressure. As mentioned above there are many technical problems which have to be settled. So, a method to use CD4 derived from human lymphocyte having specific affinity to HIV, as absorbed carrier in column for plasmapheresis is also proposed. It cannot be used as a medical procedure because of the lost affinity due to decay by autoclave treatment. In addition, there are also methods using thermostable molecules, a high molecule polymer or color ligand as an affinity carrier to HIV. However, as these molecules do not originally have specific binding ability to HIV, they cannot be used because they bind to blood ingredients faster than to HIV.\nIn this way, aiming at the development of an HIV treatment medicine, research to produce a vaccine and neutralizing antibody is flourishing, but useful medicine has not yet been developed.\nThe inventors paid attention to this present situation and developed a superior peptide to have the same degree or more affinity for gp120 compared to antibodies and to be resistant to autoclave treatment, and have already made a patent application (Japanese Patent Application No. H 8-351474 and Japanese Patent Application No. 8-351475). This peptide basically consists of a three amino acid sequence, but from a study of the sequel, we found that an affinity to gp120 of this peptide deteriorated by number and a kind of the amino acid which ranged in it. So, we knew that we needed to develop a more stable peptide."}
-{"text": "The use of lightweight composite materials in vehicles, such as aircraft, provide an improved strength-to-weight ratio that translates to fuel savings and lower operational costs. Composite materials, however, do not readily conduct away extreme electrical currents and electromagnetic forces generated by lightning strikes. Furthermore, composite or metallic structures used on vehicles are typically assembled using metallic fastener systems, which are conductive and therefore create electromagnetic effect (EME) design considerations. For example, fastener sparking modes must be designed for lightning conditions that include hot particle ejection and arcing between the fastener and surrounding structures."}
-{"text": "Conventional projection systems typically use 3 primary colours (red, green and blue) for the reproduction of colour images. The colour gamut that can be produced by an additive combination of the 3 primary colours nevertheless is limited. The colour gamut depends on the dominant wavelength (hue) and on the excitation purity (saturation) of each of the primary colours. The visible colours, the primary colours and the gamut of produced colours are usually represented in a chromaticity diagram e.g. CIE 1931 chromaticity diagram or CIE 1976 U.C.S. (Uniform Chromaticity Scale).\nOne solution to produce a wider colour gamut is by increasing the excitation purity of the primary colours, or in other words to narrow the spectral pass band of each of the three primary colours. An alternative solution to produce a wider colour gamut may be the use of more primary colours, such as e.g. using 4, 5, 6 or more colours instead of 3 colours. In printing devices, where subtractive mixing of colours is performed (by different ink cartridges), the commonly used colours are cyan, magenta, yellow and black (CMYK). In special printers, dedicated for certain types of printing, additional colours such as e.g. indigo also are used. In display technology, several techniques are known to apply more than three primary colours. One option is to generate a wider colour gamut by generating more than three primary colours by filtering them from a white illumination source and modulating them sequentially according to image data using a single chip modulation system. In such a system, typically the amount of illumination from the illumination source that is not used for displaying the image is relatively high resulting in a less efficient system. In another option, the wider colour gamut is generated by guiding more than three primary colours each to a primary colour dedicated modulator, where the primary coloured sub-beam is modulated. The latter typically results in complex and expensive systems. An example of such a system is a display system using four primary colours, each primary colour sub-beam modulated by its own primary colour dedicated modulator, as described for example in \u201cFour primary colour projection display\u201d by Roth and Caldwell in SID 05 Digest (2005) 1818.\nAdding more primary colours also introduces an electronic puzzle as to how to describe the colour space, as each colour to be displayed needs to be produced by combining the different primary colours, i.e. by subtracting in the case of printing, or by adding in the case of a display system. A method of handling the data conversion for generating images to be displayed using more than three primary colours has been described e.g. in European Patent EP 0 897 641 B1 by BARCO N. V."}
-{"text": "Patent Document 1 discloses a semiconductor manufacturing device. The semiconductor manufacturing device described in Patent Document 1 includes a stage and a plurality of pins. The stage has a plurality of holes into which the pins are respectively inserted.\nSome of the pins are used for vertically moving a substrate to be processed, which is mounted on the stage, with respect to the top surface of the stage. The other the pins are used for vertically moving a focus ring mounted on the stage with respect to the top surface of the stage.\nPatent Document 1: Japanese Patent Application Publication No. 2006-196691\nIn the above semiconductor manufacturing device, a plurality of driving units for vertically moving the pins is required as many as the number of the pins, so that the device becomes complicated. Further, in the above semiconductor manufacturing device, holes as many as the number of pins need to be formed in the stage. However, it is preferable to form a smaller number of holes in the stage."}
-{"text": "This invention relates to a propulsion unit for an inboard-outboard motor and more particularly to an improved propulsion unit that lends itself to reversal of the rotation of the propulsion device in a simple and expedient manner for twin outboard drive arrangements.\nIt is well known that a marine outboard drive generally has a side thrust due to the direction of rotation of either the propeller or other propulsion device. This side thrust may be countered by employing twin counterrotating outboard drives. For this and a variety of other reasons, it has been increasing practice to use such twin outboard drives, be they outboard motors or the outboard drive unit of an inboard-outboard arrangement.\nIn connection with such twin outboard drives, normally there is a separate internal combustion engine that powers each outboard driven. In order to achieve the counterrotation of the outboard drives, either the internal combustion engines must rotate in opposite directions or the gearing associated in the drive between the engine and the propulsion device must include a reversing mechanism for reversing the direction of rotation. In addition, each outboard drive normally includes its own forward, neutral, reverse transmission so as to permit propulsion of the watercraft in either forward or reverse directions.\nBecause of the fact that the outboard drives, either outboard motors or inboard-outboard drives, may be used either singly or in pairs, it is very desirable if the same basic construction can be utilized for both single and twin installations. This presents problems in connection with twin installations since, as aforenoted, the drives should rotate in opposite directions in such applications.\nIt is, therefore, a principal object of this invention to provide an improved and simplified arrangement for permitting reverse rotation of a marine outboard drive.\nIt is further object of this invention to provide a marine outboard drive that lends itself to ease in reversing the direction of rotation without necessitating major changes to the overall construction.\nIn conjunction with most conventional outboard drives, they are designed so that the input shaft rotates in a constant direction and the driving thrust on the unit always apply in the same direction. However, when the drive is designed so as to be driven in reverse directions for facilitating application with twin drives, the previously proposed constructions have not been completely satisfactory.\nIt is, therefore, a still further object of this invention to provide a marine outboard drive which can easily be rotated in either of two selected directions and wherein the mechanism is designed so as to take loadings regardless of the direction of drive."}
-{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent it is described in this background section, as well as aspects of the description that may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure.\nA vehicle may include an internal combustion engine and/or one or more electric motors that generate torque. A transmission selectively transfers torque to one or more wheels of the vehicle. An engine control module (ECM) controls operation of the engine. The ECM or another control module, such as a motor control module, may control an electric motor. A transmission control module (TCM) controls the transmission. The vehicle may also include one or more other control modules, such as a chassis control module, etc.\nThe control modules may communicate with one another via a network that can be referred to as a car area network (CAN). The control modules may communicate, for example, to share data. One or more of the control modules may make a decision and/or take action based on the shared data.\nThe vehicle also includes a communications module that is connected to the car area network. The communications module may also selectively establish a connection to a remote server. The communications module may update one or more of the control modules based on data from the remote server. The communications module may selectively output one or more vehicle parameters to the remote server."}
-{"text": "Databases, such as relational databases, are commonly used for managing large amounts of data, and a database user, such as a developer or end-user, can access, modify or add data within the database by way of a query. For example, in the Structured Query Language (\u201cSQL\u201d), a query may have what is referred to as multi-part identifiers in the conventional \u201cdot notation,\u201d such as \u201cA.B.C.D. \u201d As the name implies, such identifiers has are composed of multiple parts (e.g., A, B, C and D) separated by \u201cdots,\u201d or periods. Such a multi-part identifier is typically read from right to left, where the rightmost part is a column name, preceded by a table name, schema name and database name. Thus, in the above example, D would be the column name, C the table name, B the schema name and A the database name. Less than the entire path may be specified if the user is currently using the database within one of the specified levels. For example, if the user is currently operating in the table in which the desired information resides, then simply entering D as the query would result in the desired column. If the user is operating in the same database, but in a different schema, then a query of B.C.D would be required.\nAs databases become more advanced, a user is provided with greater flexibility, both in terms of the types of data that may be requested and the manner in which a query may be made to request the data. For example, User Defined Types (\u201cUDTs\u201d) are user-defined, structured data types that may be implemented using the conventional dot notation discussed above. A multi-part identifier of A.B.C may be used to query a database for a \u201cproperty C of property B of column A\u201d or a \u201cproperty C of column B of table A.\u201d These are in addition to the \u201clegacy\u201d multi-part identifier meaning discussed above.\nA database that is attempting to resolve the A.B.C multi-part identifier discussed above does not know which of the above-noted interpretations the user is intending to achieve. In most cases, only one of the interpretations will achieve a result (i.e., will \u201cbind successfully\u201d). This is because typically there will only be one interpretation that will bind for every identifier specified in a query. Thus, one interpretation of the A.B.C multi-part identifier may look for column C of table B of schema A. If one or more of the identifiers do not correspond to a value in the database, then the database will assume that this interpretation is not correct. The database may then attempt to bind an interpretation of the A.B.C query that looks for a property C of property B of column A, or for a property C of column B of table A, and so forth. If only one of the interpretations binds successfully, the database may assume that such an interpretation is the user's desired interpretation, and may return a result.\nUnfortunately, in some scenarios more than one interpretation may bind successfully. Using the above A.B.C multi-part identifier example, if the database being queried has both a column C of table B of schema A, and a property C of property B of column A, the database will not know which interpretation was desired by the user. Conventionally, this conflict was resolved by using the \u201cmost local\u201d interpretation. In other words, the interpretation that contains the least number of column prefixes would be assumed to be the correct interpretation. In the above example, such an assumption would give the highest precedence to the interpretation of \u201cproperty C of property B of column A,\u201d followed by \u201cproperty C of column B of table A,\u201d and then \u201ccolumn C of table B of schema A.\u201d Once this assumption is made, the result is returned to the user.\nSuch an algorithm has a significant shortcoming. Namely, if the assumption returns an incorrect result (e.g., property C instead of the intended column C), the user may not know that the returned result is incorrect. This shortcoming is particularly acute when the returned result is in the same format as the intended result. For example, if the returned and intended results are both numerical values, there may be no way for a user to realize that the returned result is incorrect by merely looking at the result. The incorrect result might then go on to corrupt other data if it is used in calculations or the like. Tracking down the source of such an incorrect result may prove to be extremely difficult for a user because there will be no outward sign that the multi-part identifier resulted in two possible interpretations.\nIn other scenarios, none of the possible interpretations bind successfully. Conventionally, in such a case a generic error message such as \u201ccould not find A.B.C\u201d would be issued. Such a message is typically not very helpful to a user because it does not describe the interpretation(s) that were attempted, and what the exact problem with each interpretation was.\nAccordingly, it is desirable to have a mechanism for discerning user intent when interpreting a multi-part identifier such as a UDT. In addition, it is desirable to have a mechanism that provides useful error messages to a user in the event of an interpretation conflict, failure or the like. The present invention satisfies these needs."}
-{"text": "Diisocyanates which contain urea groups are generally distinguished by a relatively high melting point. Due to the low vapor pressure thereof and due to the presence of the urea groups, they are preferably used as starting materials for the production of highgrade polyurethane/polyurea plastics. These diisocyanates are usually produced by reacting 2 mols of diisocyanate with 1 mol of water in an inert solvent. It is preferred to use diisocyanates which have two differently-reactive groups, when a definite reaction course is to be achieved. Thus, French Pat. No. 1,103,329 describes the production of bis-(3-isocyanato-4-alkoxy-, -alkyl- or chloro-phenyl) ureas by a reaction of diisocyanate and water in a molar ratio of 2:1, in a solvent, which solvent must not contain reactive hydrogen atoms (for example in the form of hydroxyl groups). Ethyl acetate, diethyl ether or acetone are preferably used as solvent.\nU.S. Pat. No. 3,906,019 describes a process for the production of monomeric di-(isocyanatotolyl)-urea, in which the reaction is carried out in an excess of one of the two non-miscible reactants (toluylene diisocyanate and, preferably, water) at from 20.degree. to 40.degree. C. The urea which is not soluble in either of the starting materials is separated in a yield of from about 60 to 85%, mainly in monomeric form, and it must then be purified. Lewis acids or Lewis bases may be used as catalysts for the reaction described therein, but pyridine is preferred\nAll these processes are carried out in solvents. The diisocyanates must be very soluble and the water which is added must be at least partly miscible in the solvents. The solvent must not exert a polymerizing effect on the isocyanate and must be free from isocyanate-reactive functional groups. The urea diisocyanates are produced as relatively insoluble compounds and are isolated by filtration. For further processing of these urea isocyanates for polyurethane production, it is necessary to convert the urea isocyanates obtained by filtration (and freed from solvents under vacuum) into a finely-divided form by a suitable grinding operation. As a result of the high melting point and the difficult solubility of these isocyanate ureas, inhomogeneous products are often obtained during the reaction.\nAccording to German Auslegeschrift No. 2,902,469, suspensions of (poly)-urea-diisocyanates in isocyanateprepolymers are produced by dissolving 1 mol of aromatic diisocyanate in the NCO-prepolymer and reacting it with from about 0.4 to 0.8 mols of water. Organic diisocyanates which have differently-reactive NCO groups are used to produce the NCO-prepolymers, as well as to produce the (poly)urea diisocyanates.\nNo process has been known in which a polyisocyanate containing urea groups is directly produced by reacting organic polyisocyanates with water in an NCO-reactive organic solvent, for example long-chain hydroxyl polyethers or hydroxy polyesters. According to the prior art, such a process would not seem very promising since the hydroxyl groups of the polyether or polyester generally have approximately the same reactivity with respect to NCO groups as water. Thus, it would be expected during the reaction of a polyisocyanate with water in a polyether or polyester for the NCO groups to react to about the same extent with the water and with the OH-groups of the polyol.\nThe separate production of urea-isocyanates, the isolation and purification thereof and the subsequent dispersion or dissolution of the urea-isocyanates in relatively high molecular weight polyols, such as are used for the synthesis of polyurethanes, is expensive. An object of the present invention is to synthesize urea isocyanates directly in the relatively high molecular weight polyols from aromatic isocyanates which are easily accessible in a simple and economic process, in order to be able to use these urea isocyanate/polyol dispersions or solutions directly for the synthesis of polyurethanes."}
-{"text": "Many vehicles, such as recreational vehicles, currently have furniture that includes a human support surface (e.g., beds, couches, chairs, etc.) that an occupant can sit and/or lay on to occupy the furniture. However, the furniture is typically affixed to one or more surfaces (e.g., the floor and/or one or more walls) of the vehicle. Where the furniture is affixed to the vehicle, the owner of the vehicle and/or an occupant of the furniture cannot relocate the furniture within the vehicle without having to detach the furniture, move the furniture to the desired position, and reattach the furniture to the vehicle.\nAdditionally, where the furniture is affixed to the vehicle, when a person is occupying the furniture, and the vehicle experiences accelerations as a result of one or more maneuvers of the vehicle, the person may shift position because of changes in their inertia, and may be disturbed as a result. This can be problematic, particularly where an occupant of the furniture is performing any sort of task that requires attention to detail and/or fine motor skills."}
-{"text": "When an optical storage device such as a Digital Versatile Disk (DVD) drive or a Compact Disk Read Only Memory (CD-ROM) drive is instructed by the host computer to retrieve data from a disk, the device will in most cases need to move an optical pickup or head to a different radial position. A servo algorithm implemented in firmware is used to command radial movement to an actuator to which the head is attached. To complete the servo loop, the current radial position of the head must be known. Information is stored on an optical disk in concentric or spiral tracks and position information is found by counting track crossings with a track count block 100 as shown in the block diagram of a typical optical storage system 10 in FIG. 1 while various waveforms are illustrated in FIG. 2.\nThe conventional optical storage device 10 includes an optical pickup and associated circuitry 12 for receiving optical signals from a disk and generating corresponding RF signals 200, a read channel 14, a decoder 16 and an interface 18 for transmitting/receiving data and command signals to/from a host device (not shown). Additionally, the device 10 includes an envelope detector 20 coupled to receive the RF signals 200 and generate a signal T-B.sub.-- ENV 202 representing the difference between the amplitude of the top of the RF signal envelope 200 and the amplitude of the bottom of the RF signal envelope 200.\nIf the optical device 10 reads DVD disks, a differential phase detector 22 generates a DVD.sub.-- PES (position error signal) 210 from phase detection signals A, B, C and D received from the pickup 12. If the optical device 10 reads CD-ROM disks (either exclusively or in addition to reading DVD disks), a differential amplifier 24 and an analog-to-digital converter 26 generate a CD-ROM.sub.-- PES 206 from phase detection signals E and F received from the pickup 12.\nExemplary waveforms of the RF signal 200, the T-B.sub.-- ENV signal 202, CD-ROM.sub.-- PES 206 and the DVD.sub.-- PES 210 are illustrated in FIG. 2. Each of these signals is input into the track count block 100. The conventional optical storage device 10 typically employs a dual-phase track counting method and requires that the T-B.sub.-- ENV signal 202 be filtered by an RF filter 101 and \"sliced\" by a slicer 102 relative to a threshold 204, resulting in an RF.sub.-- SLICE signal 212. Additionally, the appropriate PES (depending upon the type of disk in the storage device 10) must be selected by a multiplexer (MUX) 108, filtered by a PES filter 110 and sliced by a slicer 112 relative to a threshold 205 to generate a PES.sub.-- SLICE signal 214.\nCount logic 104 causes a counter 106 to increment or decrement depending on the phase relationship of RF and PES slice signals 212 and 214. When the optical pickup 12 moves in one direction relative to the tracks, the PES.sub.-- SLICE signal 214 leads the RF.sub.-- SLICE signal 212 by 90 degrees; if the pickup 12 moves in the opposite direction, the PES.sub.-- SLICE signal 214 lags the RF.sub.-- SLICE 212 by 90 degrees. (Although only the phase of the PES signal 206 changes when the track crossing direction changes, because of physical imperfections in the disk, a sinusoidal component called runout may be superimposed on the seek profile. The effect will be that the observed track count appears as if it is opposite to the intended direction. For a track counting system to be accurate, changes in direction must be accurately accounted for; consequently, both the RF and PES signals must be used.)\nWhen a seek is to be performed, servo firmware loads the counter 106 with a tracks-to-go number representing the number of tracks between the current track and the target track. The firmware also asserts a seek-in-progress signal to the counter 106 allowing the counter 106 to count down in response to a signal from the count logic module 104 and output a track count signal 216 having a resolution of 1/4 track. Additional counter input signals are seek-direction and quadrature-direction which are used to change the track count direction for a given RF to PES phase relationship. The spiral-compensation signal, also an input to the counter 106, increments or decrements the track count each time an index occurs during a seek, determined by the disk's spiral direction.\nThe dual-phase track count method has worked well in CD-ROM applications but has numerous drawbacks when applied to DVD. Maximum seek velocities in DVD devices are much higher than typical CD-ROM seek velocities and can cause unreliable detection of the RF modulation. For example, at 500,000 tracks per second in a 1.times. continuous linear velocity (CLV) DVD device, the actual velocity of the media under the head can be as slow as 80% of nominal. Assuming that there are 26,175,000 bits per second and that one-fourth of the expected RF modulation period is needed to define a peak, then ##EQU1## channel bits are available to define the peak. Because the average wavelength is 4.7 channel bits, an average of only one peak will be available to be peak detected by the envelope detector. Taking resolution effects into account, the peak detected signal is likely to be smaller than when the modulation is slower. Therefore, at high velocities, many of the track crossings will be missed by the RF.sub.-- SLICE signal 212.\nAnother drawback is that DVD signals are noisy and incompatible with dual-phase. The depth of modulation in the DVD RF signal 200 is specified to be smaller than the corresponding signal in a CD-ROM device. The signal to noise ratio is thereby reduced at the input to the RF slicer 103 and the track count error rate is increased. The purpose of the filter 10 is to reject some of the noise. Because the signal of interest is a sine wave with a widely varying frequency, an optimum filter would be a tracking bandpass filter. However, conventional technology does not currently allow a practical implementation of such a filter; rather, a low pass filter is typically used but at high seek rates, the bandwidth is greater than necessary, thereby permitting low frequency noise to pass. While the bandwidth can be changed, intervention by firmware is required.\nAnother measure of the error rate margin is the relative spacing of the edges of the slice signals 212 and 214. The rising edge of the signal 214 should be in the center of the positive pulse of the signal 212 and every edge should be in the center of an opposing signal pulse. Thus, the duty cycle of the signal should be 50% for constant seek velocities. To meet this constraint, the slicer threshold should be halfway between the maximum excursions of the signals at the input to the slicers 102 and 112. Because the depth of modulation can vary from disk to disk, the RF slicer threshold must be calibrated for each. However, even when calibrated, the threshold will be non-ideal in certain regions of the disk and the slice signals 212 and 214 will have poor duty cycles. The problem is exacerbated by the sawtooth-wave form of the DVD.sub.-- PES 210.\nTo compensate, a firmware Kalman filter can be used to estimate the track count and correct errors. However, a firmware Kalman filter can generally only detect large, defect-caused errors and can fail to detect small, noise-induced track slips. The filter also is heavily dependent on the track counts being received at regular intervals. If the T-B.sub.-- ENV or PES thresholds are incorrect, the RF or PES slice duty cycle will not be 50% and the count steps will be irregular, necessitating a larger error detection threshold. It is especially difficult for the firmware filter to detect errors as the filter instruction rate is 88,000-176,000 instruction cycles per second while the maximum track crossing rate is greater than 500,000 tracks per second.\nMoreover, current systems use no information about any previous state of the system. Consequently, errors in the track count can resemble instantaneous changes in velocity even though such changes are physically impossible."}
-{"text": "One common style of prior art poultry house is the so called post-style poultry house. These post-style poultry houses are quite prevalent and typically were built using posts prior to the wide spread use of structural trusses. A typical post-style poultry house is between about 34 and 40 feet wide and about 300 to 500 feet long. The posts in such a house are usually spaced from one another and are used to support the roof of the poultry house.\nWhile a post-style house as just described is sturdy and long-lasting, in recent years changes in poultry harvesting technology have made this less than an ideal design for continued use. For example, traditionally once the poultry are ready to be harvested from the poultry house and to be taken to a processing plant, they are gathered up and placed in cages for transport. Currently, this often is being done manually by picking up poultry and putting them in the cages. More recently, automatic harvesting equipment is being developed which constitutes a self-propelled vehicle with equipment at the front end thereof for gathering the poultry and placing them in cages. Such machinery can be obstructed in its operation by the posts. As a result, there exists a strong economic incentive for replacing the post-style houses with a clear span-style poultry house.\nUnfortunately, to simply knock down existing post-style poultry houses and replace them with new, clear span poultry houses is prohibitively expensive for many poultry farmers. Moreover, it is economically wasteful inasmuch as much of the structure of the post-style poultry house might still be in good shape and should not be demolished and discarded. Accordingly, it can be seen that a need yet remains for a method and apparatus for economically converting post-style poultry houses to clear span poultry houses. It is to the provision of such a method and apparatus that the present invention is primarily directed."}
-{"text": "The present invention generally relates to an oil pump. More specifically, the invention is an oil pump with dual scavenging for a twin cam engine.\nIt is an object of the invention to provide an oil pump with dual scavenging for a twin cam engine, with more than one pump returning oil from a sump or oil gathering location to an oil storage area or an area away from the engine.\nIt is an object of the invention to provide an oil pump with dual scavenging for a twin cam engine, with a design that improves oil control within the engine, which improves engine performance by maintaining a dry sump.\nIt is an object of the invention to provide an oil pump with dual scavenging for a twin cam engine, with a separate oil passageway from the crankcase to a first suction pump and another separate passageway from the cam chest to a second suction pump.\nWhat is really needed is an oil pump with dual scavenging for a twin cam engine, with more than one pump returning oil from a sump or oil gathering location to the oil storage area or an area away from the engine, with a design that improves oil control within the engine, which improves engine performance by maintaining a dry sump and a separate oil passageway from the crankcase to a first suction pump and another separate passageway from the cam chest to a second suction pump."}
-{"text": "Heretofore in the serving of casserole dishes or oven heated ware, there have been regular insulated pads known as hot mats which are generally regular in shape such as square or round and containing some type of an insulation material. The mat may be wrapped around the edge of a casserole to hold it when removing it from the oven and carrying it to a table. Additionally, there have in the past been gloves of an insulated type such as barbecue gloves, etc., where the casserole can be gripped for removing it from the oven.\nIn addition, when the hot casserole is placed on the table it is normally desirable to have some type of insulated material between the casserole and the table to prevent burning and discoloration thereof. This in the past has been taken care of by trivets of various types which are elevated above the table and which will allow the hot dish to sit on it without hitting the table.\nHowever, in none of the aforementioned devices is there any provision made for retaining the heat from the casserole during the time it is on the table and the food is served from it.\nIn the past, there have been devices available which do help to retain the heat in the casserole during serving such as the type where a rack is provided to set the casserole in and underneath are candles or other heating elements. Further, there have been provided in the past insulated dishes such as Thermos. Further, there are certain trivets known as electrical trivets which have a mild current passing therethrough.\nHowever, again in none of these devices is there provision for a flexible holder to grip the casserole while it is moved from the oven to the table."}
-{"text": "The art prior to the invention shown in my copending application Ser. No. 06/502,763 filed June 9, 1983, did not provide an acceptable locking/quick release socket wrench extension.\nApplication Ser. No. 06/502,763 discloses a quick release locking arrangement for a socket wrench extension having a square drive end in which a transverse bore retains lock means for movement between a lock position in which the locking means projects from one end of the bore far enough to engage a cooperating cavity or undercut in the drive recess of a wrench socket (or other tool) and a release position in which the lock means lies within the cross section of the drive end and permits removal of the drive end from the socket. A lock tab is mounted for movement between a first position obstructing movement of the lock means to its release position and to a second position in which the lock means can move to the release position. The lock tab is a part of and controlled by a sleeve which is biased to hold the lock tab in the first position.\nMy prior invention required manual actuation of the sleeve to mount a socket on the extension. The design did not have \"push-on\" mounting capability."}
-{"text": "Known hydroelectric generating equipment systems have mechanical drawbacks which substantially limit their economic utilization - particularly at low head sites. Particularly, all known types require a direct mechanical connection between the turbine and the generator. Known low-head axial turbine equipment are of three types; (1) The rim-generator type--in which the generator rotor is located on the periphery of the turbine runner; (2) the tube-type--in which the generator is located outside the water passage and connected to the turbine through an extended drive shaft; and (3) the bulb turbine type--in which both the runner and the generator are directly connected and enclosed within the water passage.\nAll three types require expensive complicated concrete Civil Works to encase and structurally support the turbine and generator in a common housing. There are also problems in matching up turbine speeds with generator speed requirements, accessibility for maintenance purposes to the important components and the inability to install most known equipment into most already existing dams or structures.\nThe following described invention proposes to solve all of the above problems and more, by allowing a physical separation of the turbine and the generator and accomplishing several other required features (such as synchronous generator speed control) through the unique application of certain fluid drive components."}
-{"text": "Digital Subscriber Line (DSL) has become popular in recent years because it can provide high-bandwidth connections between a telephone switching station and a home or office over existing telephone lines. DSL technologies use sophisticated modulation schemes to pack data onto copper wires of plain old telephone system (POTS). Among the exiting DSL technologies, Asymmetric Digital Subscriber Line (ADSL) is the most practical since it fits average users' need for greater downstream bandwidth than upstream bandwidth.\nDiscrete Multi-tone (DMT) is a method of modulating a DSL signal for transmission on a usable frequency spectrum divided into a plurality of sub-carriers or frequency channels. For example, a spectrum may be divided into 256 channels (\u201cbins\u201d) of 4.3125 kHz each. The center frequency of a bin is 4.3125 kHz multiplied by the bin number. Thus, a specific bin may be referred to by its bin number instead of its actual frequency range. Dividing the frequency spectrum into multiple channels reduces cross-talk in signal transmissions. In a DMT system, a usable frequency spectrum is typically allocated to the upstream and downstream transmissions based on a standard scheme. For example, according to International Telecommunication Union (ITU) standard G.992.1 (Annex A), the downstream transmission may occupy bins 6 through 255 and the upstream transmission may occupy bins 6 through 31. In a common operation mode, to reduce near-end cross-talk (NEXT), the downstream transmission may start from bin 33. Other standards or schemes also exist that allow an overlap in bin allocations to the upstream and downstream traffic.\nHowever, none of the fixed bin allocation schemes can consistently achieve optimal data rates. Due to the ever-changing operation environment (e.g., noise level, loop attenuation and echo rejection rate), quality of a connection based on a particular scheme usually drifts over time. Especially when the cross-talk noises become highly asymmetric, it may be difficult to rely on a fixed bin allocation scheme to maintain link performance. As a result, the available bandwidth is not fully utilized and customers may become dissatisfied with inconsistent link performance.\nIn view of the foregoing, it would be desirable to provide an efficient and cost effective solution for bin allocation that overcomes the above-described inadequacies and shortcomings."}
-{"text": "Field\nThe present disclosure relates generally to a method and apparatus for detecting license plate information from an image of a license plate and more specifically, detecting license plate information from an optical image, captured by a mobile apparatus, that includes a license plate image and several other object images.\nBackground\nIn recent years, collecting still images of license plates has become a common tool used by authorities to catch the drivers of vehicles that may engage in improper or unlawful activity. For example, law enforcement authorities have set up stationary traffic cameras to photograph the license plates of vehicles that may be traveling above a posted speed limit at a specific portion of a road or vehicles that drive through red lights. Toll booth operators also commonly use such stationary cameras to photograph vehicles that may pass through a toll booth without paying the required toll. However, all of these scenarios have a common thread. The camera must be manually installed and configured such that it will always photograph the vehicle's license plate at a specific angle and when the vehicle is in a specific location. Any unexpected modifications, such as a shift in angle or location of the camera would render the camera incapable of properly collecting license plate images.\nAdditionally, camera equipped mobile apparatuses (e.g., smartphones) have become increasingly prevalent in today's society. Mobile apparatuses are frequently used to capture optical images and for many users serve as a replacement for a simple digital camera because the camera equipped mobile apparatus provides an image that is often as good as those produced by simple digital cameras and can easily be transmitted (shared) over a network.\nThe positioning constraints put on the traffic cameras make it difficult to take images of license plates from different angles and distances and still achieve an accurate reading. Therefore, it would be difficult to scale the same license plate image capture process performed by law enforcement authorities to mobile apparatuses. In other words, it is difficult to derive license plate information from an image of a license plate taken from a mobile image capture apparatus at a variety of angles, distances, ambient conditions, mobile apparatus motion, and when other object images are also in the image, which hinders a user's ability to easily gather valuable information about specific vehicles when engaging in a number of different vehicle related activities such as buying and selling vehicles, insuring vehicles, and obtaining financing for vehicles."}
-{"text": "It is often necessary for certain telephone-specific or switching system-specific operating steps to be performed in order to program a function key. The allocation of a programmed speed dialing key to a telephone number is usually indicated here by the user manually writing on a paper label strip of the telephone with a pen.\nThere are also modern telephones whose keys are assigned an LCD or LED display (LED: Light-Emitting Diode; LCD: Liquid Crystal Display) in each case, on which the stored telephone number of the programmed speed dialing key can be displayed as an electronic key label. Manual labeling of the keys is not required for such telephones. Function keys of this type with a display are often referred to as self-labeling keys.\nIn known communication systems, setting the display for a key is then carried out during programming of the function key. For instance, following programming of a speed dialing key, the telephone number allocated to the speed dialing key is displayed on the display for the speed dialing key. In addition, as part of programming the speed dialing keys, the user is often able to manually edit the electronic key label via a user interface, by entering a text or the name of the subscriber for the stored destination telephone number for example. This name can then be displayed on the display of the speed dialing key instead of the telephone number."}
-{"text": "Many different types of toasters have been manufactured and used all over the world for decades, ranging from small single or double slice home toasters to large commercial toasters. Toasters may be used to toast sliced bread, rolls, bagels, biscuits, and other similar food products.\nOne type of toaster is a conveyor toaster that continuously moves food prod products through the appliance. Generally, these conveyor toasters include a housing having a toasting cavity, a conveyor for moving the bread/toast through the toasting cavity, and one or more heating elements.\nVarious efforts have been made to develop conveyor toasters, and the following references show a few examples. These references are incorporated by reference herein, in their entireties:\nU.S. Pat. No. 1,696,613: Toaster\nA toaster includes a horizontal conveyor, the bread to be toasted being-placed on the conveyor at the front of the toaster and carried toward the rear in close proximity to suitable heating means, whereby the bread is toasted. At the rear of the toaster the bread is fed by the conveyor to means for returning it to the front of the toaster, this means being shown as an inclined slide. The toaster further includes means for regulating a draft of cooling air around the heating means whereby the generation of heat may be-regulated for different toasting operations and for light or heavy loads, and also a special crumb tray arrangement below the heating means.\nU.S. Pat. No. 1,708,522: Toaster\nA toaster is provided in which the articles to be toasted, such as slices of bread, are fed through the apparatus by suitable conveying means, and wherein the heating of the heating elements is started by insertion of the material to be toasted. The toaster includes conveying means for the material to be toasted and means for driving or stopping said conveyer in accordance with the presence or absence of the toaster of material to be toasted, and further includes a casing having insertion openings at different positions along the conveyer so that the length of path traveled by the bread in passing the heating means may be varied to vary.\nU.S. Pat. No. 3,517,605: Bun Toaster\nA compact and easily cleaned bun toaster such as for the heels and crowns of hamburger buns and the like. A conveyor confronts a toasting platen for driving buns along the surface of the conveyor to toast them. The conveyor is removable as a unit from the frame mounting the platen, thereby to expose both the conveyor and the platen for easy cleaning. The platen is movable toward and away from the conveyor to adjust the space therebetween. The platen is stepped to facilitate simultaneous toasting of bun heels and crowns of different thicknesses. The conveyor is in driven engagement only when moved into suspended engagement with the frame.\nU.S. Pat. No. 3,693,536: Apparatus and Method of Toasting Bread Like Articles\nAn electric toasting apparatus comprised of a conveyor adapted to support articles to be toasted and having the transverse width and a preselected rate of movement in a longitudinal direction. The apparatus includes at least two radiant heater means disposed adjacent the conveyor, each arranged to provide generally uniform heating conditions across substantially the entire transverse width of the conveyor and along at least a substantial portion of the longitudinal length of the conveyor. One heater runs constantly and has a heat output to bring the articles close to the toasting point. The other is modulated to control its heat output and thus the total output of both heater means to give the desired degree of toasting. The method comprises the steps of: (a) conveying the articles past and adjacent to the first and second heater means at a uniform rate; (b) fully energizing the first heater means in an amount sufficient to bring the articles at least close to the toasting point; and (c) simultaneously energizing the second heater means to produce less than its maximum heat output, but sufficiently to fully toast the articles.\nU.S. Pat. No. 3,835,760: Apparatus for Toasting Bread-Like Articles\nApparatus for toasting bread-like articles is disclosed comprising a housing providing a toasting chamber in which electrical heating elements are disposed together with an endless conveyor operable to receive and carry articles to be toasted past the heating elements. The housing has an entrance opening for introducing articles onto the conveyor and when the articles have been toasted they are discharged from the conveyor for removal from the apparatus. The housing is provided with a discharge opening in the top wall thereof and adjacent the rear wall of the housing, and a flue arrangement is provided within the housing to induce hot air in the housing to flow through the discharge opening. Flow of hot air through the flue induces ambient air to flow into the toasting chamber through the product entrance, beneath the conveyor and heating elements and thence upwardly toward the top of the chamber. A temperature sensing element is disposed in the housing in the path of ambient air flow and is responsive to the temperature sensed to cause energization and de-energization of the heating elements.\nU.S. Pat. No. 4,286,509: Toasting Apparatus\nAn economical energy saving toasting apparatus includes a housing with an interior baffle defining a toasting chamber therein and an inlet for the introduction of the product to be toasted into the chamber. The housing also includes an outlet through which the toasted product is dispensed. A conveying assembly for conveying the product from the inlet through the toasting chamber to the outlet is also included within the housing. Mounted within the toasting chamber is a first set of resistive heating elements that are electrically connected in a circuit and continuously energized during the operation of the heating apparatus. A second set of heating elements is also included in the heating chamber and connected to a timer in order to be energized for only selected periods of time. A micro-switch is also included in the apparatus to actuate the timer and a micro-switch actuator is positioned in the inlet and adjacent to the conveying device so as to be engaged by the product prior to entry into the toasting chamber. The engagement of the product actuates the micro-switch to energize the timer and the second set of heating coils for a predetermined period of time.\nU.S. Pat. No. 5,673,610: Apparatus for Conveyorized Toasting of Breads and Like Food Items\nAn apparatus for conveyorized toasting of sandwich buns and like bread and food items simultaneously on opposite sides comprises a central heated platen with two bun transport conveyors traveling in spaced relation along opposite sides of the platen and with a pair of auxiliary heating elements disposed outwardly of the respective food transport runs of the conveyors in facing relation to the opposite sides of the platen. Each conveyor is biased toward the platen by a pair of parallel pivot arms urged by springs into pivoted engagement against the transport run of the respective conveyor to define a predetermined desirable spacing to the facing side of the platen.\nU.S. Pat. No. 5,746,116: Rapid Toasting Apparatus\nAn improved toasting apparatus is disclosed including a variable speed conveyor belt for conveying products through the toaster; a bottom heating element positioned below the conveyor belt, the bottom heating element including variable control for controlling the amount of heat provided by the bottom heating element; a drying zone including a first plurality of upper heating elements positioned above the conveyor; a toasting zone including a second plurality of upper heating elements; and wherein the bottom heating element and the first and second plurality of upper heating elements are controlled individually and independently.\nU.S. Pat. No. 5,983,785: Contact Toaster with Infinite Adjustment\nA contact toaster has a housing with a product inlet, product outlet and at least one toasting chamber defined in the housing. A heated platen is mounted in the housing and has a platen surface arranged for toasting the food products. A flexible endless belt, which when rotated, is arranged to pass the food products in a toasting path along the platen between the belt and the platen. An infinite controller acts to selectively displace the belt to adjust the distance between the belt and the platen to any value between minimum and maximum distance limits to accommodate food products of different widths. A guide is situated between the belt and the infinite controller so that the controller is in contact with at least one area of the guide. The infinite controller includes a variable distance element with a cam surface which is settable to any selected point of the cam surface between minimum and maximum cam radial limits to move the guide and belt to any belt to a selected platen distance that can have any value equal to or intermediate the maximum and minimum distance limits.\nU.S. Pat. No. 6,019,030: Conveyor-Type Food Product Toaster\nA toaster for bread-type food products includes a housing defining at least one toasting cavity. A toasting heater is operatively associated with the cavity. An endless chain-type conveyor is provided for conveying food products through the cavity. An endless belt is disposed about the conveyor to prevent the chain-type conveyor from imprinting the food products. The belt includes a seam, and a retention portion of the belt at the seam is engaged with the conveyor to prevent slippage between the belt and the conveyor.\nUS Publication Number 20100139497A1: Food Heating Device\nA food heating device includes a first conveyor having a first thermally-conductive belt and a second conveyor having a second thermally-conductive belt. A first platen is disposed on a first side of the first thermally-conductive belt and a second platen disposed on a first side of the second thermally-conductive belt. The first conveyor and the second conveyor are arranged and spaced to transport a food product between a second side of the first thermally-conductive belt and a second side of the second thermally-conductive belt.\nUS Publication Number 201000275789A1: Toaster with Removable and Adjustable Conveyors\nA dual-sided, conveyor toaster provides operator-adjustable conveyors that are also operator removable. The removable conveyor assemblies are chain driven and removably supported in the toaster by re-positionable mounting mechanisms embodied as either adjustable pins that engage fixed slots or adjustable slots that engage fixed pins. Fixed pins can be located on the conveyor and engage adjustable slots in the toaster; adjustable pins can be located on the toaster and engage fixed slots on the conveyor. The conveyor assemblies use non-marring plates to urge food products against the heated platen surfaces and to carry the food products across the platen."}
-{"text": "Commercial hardening techniques for ferrous metal components typically involve the steps of heating the components in a furnace or a bath of molten salt, such as a chloride, quenching the components in a bath of another molten salt such as a nitrate, nitrite, or mixture thereof, and thereafter causing the parts to be passed through one or more rinse baths to clean off salts which cling to the parts as a result of the heating and quenching steps. A certain amount of each of the molten salts is transferred by drag-out to the rinse bath, thus requiring occasional replenishment of the molten salts, as well as replacement of the rinse water.\nIt is desirable to recover the salts from the rinse water and return them to the heat treating baths. However, a number of problems have heretofore prevented such recycling on a straightforward and economical basis. First, it is difficult to justify economically the installation of recovery apparatus unless the quantities of reclaimed salts are very high, and the recovery process can be carried out without shut-down of the heat treating system. Secondly, it has been found that the impurities which typically find their way into heat treating baths tend to build up to unacceptable levels particularly fast when the salts are recovered from the rinse and recycled into the heat treat bath. For example, chloride salts from a heating bath are dragged into the quench bath and thence into the rinse where recycling tends to return them to the quench bath. Recycled chloride impurities quickly reach insolubility levels in the quench, and crystallize out as a sludge which must be removed from time to time. In addition, chloride tends to substantially reduce the quench efficiency of a nitrate/nitrite bath and, thus, removing the chlorides is highly desirable.\nAnother impurity which will build up quickly as a result of salt recycling is carbonate, a reaction product of decomposition of the salt itself. For example, a sodium nitrate/nitrite salt tends to form Na.sub.2 CO.sub.3, the carbon being taken out of the CO.sub.2 present in the air. The sodium carbonate is quite soluble in the rinse water but only sparingly soluble in the nitrate/nitrite quench bath and, thus forms a precipitate in the bath which impairs quenching efficiency and coats mechanical components in contact with the bath.\nIt should be understood that both the chloride and carbonate impurities, as well as many other impurities such as hard water ions, are typically present in the heat treating salt baths of prior art systems. However, the levels of concentration of such impurities are kept reasonably and acceptably low for longer periods of use by reason of the fact that impurities as well as desirable salts are continuously lost by drag-out. Lost salt, of course, is replaced by fresh salt which does not contain the impurities. Where salt is recycled, however, undesirable impurities dragged into the rinse and which are soluble in the rinse water are retained in the recycled materials and, after a short time, reach unacceptable concentrations."}
-{"text": "Fabric care, particularly cleaning, can be divided into two basic theoretical approaches, namely, water-based cleaning or dry/powder cleaning. Dry cleaning can be perceived by customers as raising certain concerns regarding solvent residue and volatile organic content (VOC), however, water-based cleaning also has disadvantages.\nFor one example, in a conventional water cleaning process, significant amounts of clean water are required for both washing and rinsing cycles. Typically, the weight ratio of consumed water to fabric is around 12:1 to 20:1. As can be appreciated, this is a challenge for clean water supply infrastructure, especially those facing water shortage issues. Moreover, many forecast water-saving as a mega-trend in sustainability. Accordingly, consumers desire to conserve water is likely to increase over time.\nThus, it is desirable to provide a fabric care composition that achieves excellent washing performance with water, whilst requiring less water than conventional water-based cleaning processes. It is also desirable to provide an efficient and environmentally friendly cleaning method compared to conventional cleaning processes."}
-{"text": "In means of transportation, in particular airplanes, storable tables are frequently provided for the passengers. In VIP airplanes and in the first class section of conventional commercial passenger airplanes, a table of this type is usually pulled out of what is known as a credenza, which is arranged next to the seats, and is folded into a use position, in which the table top is arranged partially over the seat surface of the seats.\nIt is problematical in tables of this type that a plurality of requirements have to be united which are difficult to combine with one another. The credenza is usually lower than the armrests of the seats. This means that the table top has to be arranged considerably above the upper edge of the credenza in the use position. A mechanism therefore has to be provided which holds said table top in this higher position. The table top has to be self-supporting, and all loads which act on the table top therefore have to be conducted into the credenza. This can be problematical, in particular, in a long lever arm for the active forces (for example, a passenger who stumbles along the aisle).\nIt is known to secure the table top of a table of this type pivotably on a carrying arm which protrudes out of the upper edge of the credenza in the use position and therefore produces the necessary spacing between the upper edge of the credenza and the table top. The lid of the credenza is configured as a type of catching hook, and that edge of the table top which points away from the passenger seat can be hooked into the lid in its open position, with the result that the tilting moment which acts on the table top is conducted into the credenza via this lid. This solution is disadvantageous, as a lid which is fastened to the credenza by means of a hinge is not particularly well suited for absorbing loads."}
-{"text": "Integrated semiconductor circuits (ICs) can be damaged by transient pulses or overvoltages, which are coupled in via terminals (pads) or directly in lines, such that they become nonfunctional or are even destroyed. Such pulses or overvoltages may occur for example in the event of so-called electrostatic discharges (ESD). High voltages and high currents, associated with transient or ESD disturbances, cause high disturbing powers to occur.\nSuch a pulse (e.g. burst) may also occur in many fields of application, e.g. automotive engineering. In automotive engineering, by way of example, there is the requirement that circuits of this type which have to function in the high-voltage range up to 90 volts or above are also to be designed for significantly higher disturbance pulse levels.\nFor high-voltage applications produced by means of high-voltage processes, provision is usually made of protection devices that are initiated or triggered by an electrical breakdown. The breakdown voltages must be significantly greater than the maximum permissible operating voltages of the application circuit to be protected. Only then is it possible to guarantee an undisturbed functionality of the integrated circuit. In the case of a fault, e.g. when an impermissibly high voltage is present, this overvoltage is dissipated to reference potential or ground by the protection circuit and downstream assemblies are thus protected from the high voltage.\nOne alternative for such breakdown-based protection concepts is an active circuit for protecting an integrated circuit, comprising a combination of an active trigger circuit with a known protection device such as a thyristor or a bipolar or MOS protection transistor.\nActive circuits for protecting the IC are often triggered by the rise of the transient signal. In this case, the signal rise per unit time is detected and a protection transistor or a protection circuit is turned on by means of a drive circuit.\nIn the case of a fault, the protection circuit may accordingly be understood as an actively triggered overvoltage or overcurrent surge arrestor. A fast activation of the protection circuit is necessary in the case of a fault.\nShort switch-on times and a precise switch-on threshold of the protection circuit for the integrated circuit and the protective effect thereof for different forms of disturbance pulses are important aspects of the product specification and constitute a competitive advantage.\nU.S. Pat. No. 5,982,601 discloses a thyristor (SCR\u2014silicon control rectifier) for ESD protection, which is triggered directly by the transient signal. The thyristor is realized in the semiconductor arrangement in a manner known per se by means of an n-type well, a p-type well and highly doped n-type and p-type regions. The transient voltage is detected by means of an RC element. By means of inverters connected downstream, the voltage level detected at the capacitance is converted into a control signal which drives the base of the pap transistor of the thyristor structure. As soon as the output current of the then active pnp transistor generates a sufficiently large voltage drop at a resistor, the npn transistor of the thyristor structure turns on, so that the transient pulse is dissipated from the pad potential of the I/O pin to reference potential through the low-impedance thyristor path. The thyristor then automatically remains turned on until its current falls below the holding current and the turn-off condition is met."}
-{"text": "In a conveying system for conveying small articles such as cassettes or pallets, it is frequently necessary to remove the articles from the conveyor for processing. Traditionally, an operator has manually gathered the articles as they pass along the conveyor, and manually stacked the articles in preparation for transporting to the processing equipment. The manual gathering and stacking of the articles requires considerable labor and is a substantial cost factor.\nAutomatic stacking mechanisms have been proposed in the past, such as shown in U.S. Pat. No. 3,470,996. In the stacking mechanism, as shown in the aforementioned patent, articles on a conveyor are pushed laterally by a pusher unit from a conveyor to a position beneath a stacking chute. A lift cylinder then operates to elevate the article upwardly into the chute and the elevated article is then held in the chute by a second plunger which is located opposite the pusher assembly. While devices of the type shown in the aforementioned patent can serve to stack articles one on top of each other, the mechanism requires that the articles slide against each other during movement from the conveyor to the stacking chute and this can cause substantial damage and scratching to the articles, if they are of a delicate nature."}
-{"text": "1. Field of the Invention\nThe present invention relates to a keycap, and in particular to a light-guiding method and a light-guiding assembly of a backlight source within the keycap.\n2. Description of Prior Art\nBack light source can be widely used in a liquid crystal display, a signboard or an advertisement board. In addition, it can be also applied to electronic products, such as mobile phone or personal digital assistant. Since the back light module is mounted within the keycap, the user can still smoothly operate the keycaps even in the night or without the illumination of a light source.\nThe conventional keycap is provided with a light-guiding plate in the interior of the key pad assembly, and the periphery or the bottom of the light-guiding plate is provided with a light-emitting diode. When the electronic product is in use, the light-emitting diode is lighted up, and the light is introduced into the light-guiding plate. At this time, with the light-guiding plate, the light is directed to the keycap above the light-guiding plate, so that the surface of the keycap has a light-transmitting effect. In this way, the user can clearly see the character, numeral or symbol on the surface of each keycap.\nRecently, with the development of multi-color light-emitting diodes, they are also widely used in the keycaps of electronic products. The multi-color light-emitting diode generates the lights with various colors under the control of the circuit of the electronic product to, so that the surface of the keycap or some keycap groups can represent a colored back light effect with different colors. Since there is only one layer of light-guiding plate, if several multi-color light-emitting diodes are lighted up at the same time, and the lights with various colors are simultaneously introduced into the light-guiding plate, these lights may be mutually interfered with one another or mixed easily. As a result, the light displayed on the surface of the keycap is not the originally preset light. Therefore, in order to solve the mutual interference or mixing of the lights, adhesion of black patches or coating of inks is applied to the light-guiding plate so as to shield the light. After several multi-color lights are introduced into the light-guiding plate, the surface of the keycap or some keycap groups may represent a backlight effect with various colors and do not generate the interference or mixing of the lights. However, if the adhesion or coating process is not perfectly performed, such light-shielding measure may cause the leakage of light and in turn the waste of time and process in manufacturing the light-guiding plate."}
-{"text": "The present invention relates to the filtration field, and more particularly, to an improved centrifugal filtration apparatus for filtering and concentrating a solution. The improvement comprises eliminating any retentate pockets below the filter and instead employing a port on the filtrate side of the filter that is above the bottom level of the filter, thereby not permitting filtering to dryness.\nThe filtration of fluids may be accomplished through the use of filtration devices which utilize microporous filters to filter and concentrate a macro-molecular solution is well known. This technique has been used in centrifugal filtration apparatuses that rely on centrifugal forces to create pressure in the apparatus to force solutions through a filter which separates liquid solutions into filtrate and concentrate.\nThere are certain drawbacks, however, associated with conventional centrifugal filtration apparatus. While many apparatus designs can prevent filtering to dryness in certain applications, they do not in other applications, such as a swinging bucket centrifuge. Those that work in both a fixed angle centrifuge and a swinging bucket centrifuge without filtering to dryness may present the problem of either retaining a different amount of concentrate when used in a fixed angle centrifuge than they will when used in a swinging bucket centrifuge or retaining a different amount of concentrate when used in fixed angle centrifuges with different fixed angles.\nSome of the conventional centrifugal filtration apparatus require a second spin to remove the retained concentrate (i.e. dead stop volume). In these apparatus it is difficult or impossible to remove the retained concentrate from the concentrate tube with a pipette.\nAnother problem with some of the conventional centrifugal filtration apparatus is that as the concentrate volume approaches its final retained volume, the active filter area approaches zero. Therefore, the filtration rate will slow dramatically as the concentrate volume approaches its final retained volume.\nAnother problem with the conventional centrifugal filtration apparatus is that they are open systems (i.e. they contain a venting means that vents to the atmosphere in the centrifuge).\nYet another problem with some of the conventional centrifugal filtration apparatus is that they are not scaleable (i.e. they are designed to be used as small volume centrifugal filters, or large volume centrifugal filters, not both).\nCertain types of filtration devices, such as that disclosed in U.S. Pat. No. 4,632,761 to Bowers et alia, are capable of preventing filtration to dryness and contain a dead stop feature which causes filtration to cease while there is concentrate remaining within the apparatus. This device, however, filters to dryness when spun at a 90xc2x0 angle and therefore the dead stop feature will not work if the device is spun in a swinging bucket centrifuge. Also, in this type of filtration device, the amount of concentrate remaining after dead stop is dependent upon the angle of the centrifuge rotor. The filter surface area in this type of device is limited by the diameter of the device, and the surface area is relatively small when compared to the volume of liquid solution within the housing. Another problem with this type centrifugal filtration apparatus is that is an open system (i.e. it contains a venting means that vents to the atmosphere in the centrifuge). This type of filtration device is conductive to clogging because the heaviest and denser molecules within the liquid solution are forced into the membrane filter. Accordingly, this device is limited because it will filter to dryness when spun in a swinging bucket centrifuge, it will filter to different dead stop volumes when used in fixed angle rotors with different angles, it is an open system and will vent potentially harmful gases to atmosphere during filtration, and its active filter surface area is limited by the diameter of the device.\nThe types of filtration devices disclosed in U.S. Pat. No. 4,722,792 to Miyagi et alia are capable of preventing filtration to dryness and contain a dead stop feature which causes filtration to cease while there is concentrate remaining within the apparatus. In this type of filtration device, the amount of concentrate remaining after dead stop is dependent upon the angle of the centrifuge rotor. Therefore, the dead stop volume will be different when the device is used in a swinging bucket rotor centrifuge than it will be when used in a fixed angle rotor, and will also be different when the device is used in fixed angle rotor centrifuges of different rotor angle. This means that the results obtained from this type of device when used in one type of centrifuge rotor cannot be compared to the results obtained from this type of device when used in another type of centrifuge rotor. Another problem with this type centrifugal filtration apparatus is that is an open system (i.e. it must contain a venting means that vents to the atmosphere in the centrifuge). Another problem with this type of centrifugal filtration apparatus is that the filtration rate starts out high because of its relatively large filter surface area. However, as the concentrate volume approaches its dead stop volume, the active filter area approaches zero. Therefore, the filtration rate will slow dramatically as the concentrate volume approaches its final retained volume. Accordingly, this device is limited because it will filter to different dead stop volumes when used in centrifuge rotors with different angles, it is an open system and will vent potentially harmful gases to atmosphere during filtration, and the filtration rate will slow dramatically as the concentrate volume approaches the dead stop volume.\nThe types of filtration devices disclosed in U.S. Pat. No. 5,112,484 to Zuk are capable of preventing filtration to dryness and contain a dead stop feature which causes filtration to cease while there is concentrate remaining within the apparatus. In this type of filtration device, the amount of concentrate remaining after dead stop is dependent upon the angle of the centrifuge rotor. Therefore, the dead stop volume will be different when the device is used in a swinging bucket rotor centrifuge than it will be when used in a fixed angle rotor, and will also be different when the device is used in fixed angle rotor centrifuges of different rotor angle. This means that the results obtained from this type of device when used in one type of centrifuge rotor cannot be compared to the results obtained from this type of device when used in another type of centrifuge rotor. Another problem with this type centrifugal filtration apparatus is that is an open system (i.e. it must contain a venting means that vents to the atmosphere in the centrifuge). Another problem with this type of centrifugal filtration apparatus is that as the concentrate volume approaches its final retained volume, the active filter area approaches zero. Therefore, the filtration rate will slow dramatically as the concentrate volume approaches its final retained volume. Although the retained concentrate can be removed from this device with a pipette, the retained concentrate is best removed from this type of device using the concentrate cup with a second spin. This type of device is not well suited to filter small volumes. Accordingly, this device is limited because it will filter to different dead stop volumes when used in centrifuge rotors with different angles, it is an open system and will vent potentially harmful gases to atmosphere during filtration, the filtration rate will slow dramatically as the concentrate volume approaches the dead stop volume, it is not designed to filter small volumes, and it is not easy to remove the dead stop volume with a pipette.\nThe types of filtration devices disclosed in U.S. Pat. No. 5,490,927 to Herczeg are capable of preventing filtration to dryness and contain a dead stop feature which causes filtration to cease while there is concentrate remaining within the apparatus. In this type of filtration device, the amount of concentrate remaining after dead stop is dependent upon the angle of the centrifuge rotor. Therefore, the dead stop volume will be different when the device is used in a swinging bucket rotor centrifuge than it will be when used in a fixed angle rotor, and will also be different when the device is used in fixed angle rotor centrifuges of different rotor angle. This means that the results obtained from this type of device when used in one type of centrifuge rotor cannot be compared to the results obtained from this type of device when used in another type of centrifuge rotor. Another problem with this type centrifugal filtration apparatus is that is an open system (i.e. it must contain a venting means that vents to the atmosphere in the centrifuge). This type of device is not well suited to filter small volumes. This type of device utilizes the entire filter surface area until the dead stop volume is reached when used in a swinging bucket centrifuge. However when this type of device is used in a fixed angle rotor centrifuge the active filter area will approach zero as the concentrate volume approaches its final retained volume. Therefore, the filtration rate will slow dramatically as the concentrate volume approaches its final retained volume when used in a fixed angle centrifuge. This type of device requires a second spin to remove the retained concentrate solution. Accordingly, this device is limited because it will filter to different dead stop volumes when used in centrifuge rotors with different angles, it is an open system and will vent potentially harmful gases to atmosphere during filtration, the filtration rate will slow dramatically as the concentrate volume approaches the dead stop volume when used in a fixed angle rotor centrifuge, it is not designed to filter small volumes, and it requires a second spin to remove the retained concentrate solution.\nFinally, the art discloses the use of a lower retentate pocket to address the issue of filtering to dryness with centrifugal filtration devices. U.S. Pat. No. 5,647,990 to Vassarotti reveals a device having a concentrate retention pocket located at the bottom of a pre-filtration chamber, below the level of the filter membrane. Fluid flows through the filter into filtrate outlet channels, exiting these channels into the filtrate reservoir. These exit channels do not collect filtrate.\nIt is therefore an object of the present invention to provide a filtration device that may be used in a swinging bucket centrifuge, as well as a fixed angle centrifuge, without filtering to dryness.\nIt is also an object of the present invention to provide a filtration device that will filter to approximately the same dead stop volume when used in either a swinging bucket rotor centrifuge, or when used in a fixed angle rotor centrifuge.\nIt is also an object of the present invention to provide a filtration device having a relatively high filtration membrane surface area thereby enabling filtration to occur at a higher rate.\nIt is also an object of the present invention to provide a filtration device that continues to use a relatively high filtration membrane surface area until the dead stop volume is attained, thereby enabling filtration to occur at a higher rate throughout the entire filtration process.\nIt is also an object of the present invention to provide a filtration device, which minimizes the clogging of the semi-permeable membrane thereby maximizing filter efficiency.\nIt is also an object of the present invention to provide a filtration device from which the retained concentrate solution can be easily removed with a pipette.\nIt is also an object of the present invention to provide a filtration device which can be used as a closed system (i.e. does not vent to atmosphere), or that can be used as an open system (i.e. vents to atmosphere).\nIt is a further object of the present invention to provide a filtration device, which can be manufactured as a disposable small volume device, or as a disposable large volume device.\nThe foregoing problems of the prior art are solved, and the objects of the present invention are achieved, by use of a filtration apparatus constructed in accordance with the principles of the present invention.\nIn accordance with the present invention, the filtration apparatus for filtering and concentrating a solution includes a concentrate tube, which is divided into two sections. The upper cylindrical section acts primarily as a reservoir for unfiltered solution and directs unfiltered solution from the upper section into the lower section of the concentrate tube. The upper section is shaped so that all of the solution in the upper section will flow into the lower section regardless of the centrifuge rotor angle. The lower section of the concentrate tube is mainly rectangular, but may be triangular shaped at the bottom. The volume of the lower section is made as small as possible to insure that the entire filter surface area is used to filter the maximum amount of solution. The lower section may contain a pipette tip channel to facilitate the removal of the retained concentrate with a pipette. The lower section is open on one side, or two sides.\nA micro-porous membrane or membranes may be sealed to the lower section in several ways. First, the membrane can cover the open side of the lower section of the concentrate tube, or the membranes can cover the open sides of the lower section of the concentrate tube. Alternately the micro-porous membrane may be sealed to the filter cover, or covers. A second alternative is to seal the micro-porous membrane or membranes between the filter cover and a filter sealing gasket or gaskets. A third alternative is to seal the micro-porous membrane between the filter cover and the concentrate tube using a compression rib that compresses a portion of the outer periphery of the micro-porous filter between the filter cover and the concentrate tube.\nA filter cover is sealed to the open side of the lower portion of the concentrate tube, or two filter covers are sealed to the two open sides of the lower portion of the concentrate tube. The filter cover or covers contain a filter support means such as filter support ribs. The filter cover or covers also contains a filtrate outlet port. The position of the filtrate outlet port or ports determines the retained concentrate volume. The concentrate tube may contain a port or ports to allow displaced filtrate air to flow into the upper part of the concentrate tube to replace the filtered solution that flows from the concentrate tube into the filtrate tube.\nThe means for collecting the filtrate may include a filtrate tube. The concentrate tube may be sealed into the filtrate tube. The filtrate tube cap may contain a means to nest the concentrate storage tube."}
-{"text": "This invention relates in general to bird feeders and animal feeders, and is more particularly directed to demand feeders for agricultural birds, e.g., fowl such as chickens, ducks, turkeys, peafowl, or many other exotic or native birds, or small animals like cats and rabbits.\nThe invention is more particularly concerned with a feeder that dispenses feed on demand from a raised hopper onto a feed tray or catch dish that is suspended beneath the hopper, where the dispensing occurs by action of the bird or birds feeding, i.e., by action of the birds pecking and grabbing at food on the feed tray.\nIt is a significant object of the invention to create a bird feeder that is simple, easy to use, refill, and clean, and which allows for choice of a wide variety of feed types, from small grain, such as millet seed or mash, to larger particles such as pellets, corn, or dog food.\nThe bird feeder provides on-demand gravity feed and engages the birds in pecking and grabbing activities which helps stimulate normal, healthy behavior. The on-demand feeder also reduces over-feeding and eliminates food waste, dispensing feed only when needed."}
-{"text": "Recently, digital discs have been widely used in many fields. For example, a CD (abbreviation of Compact Disc) is used for digital audio records. A CD-ROM (ROM; abbreviation of Read Only Memory) is used for data memories. An LV (abbreviation of Laser Vision) disc is used for video records. In those digital discs, recorded data or information are reproduced at a speed specified to each kind of the digital discs. Thus, reproduction apparatus of digital discs are designed to reproduce the digital discs at the sole reproduction speed, i.e., a standard speed specified to each of the digital discs.\nReferring now to FIG. 1, an example is shown of the conventional reproduction system for digital discs, e.g., a CD dubbing system. The CD reproduction system is provided for carrying out a dubbing from CD players to analog tape recorders.\nIn FIG. 1, the CD reproduction system comprises a CD player 10 and an analog tape recorder, e.g., a cassette tape recorder 11. The CD player 10 includes a CD reproducing mechanism section 12, a clock generator 13 and a signal processor 14.\nThe CD reproducing mechanism section 12 comprises a disc drive motor 15, a disc turntable 16 and a pickup device 17. The clock generator 13 generates a clock pulse Pck with a prescribed frequency Fck. The clock pulse Pck is applied to the signal processor 14.\nThe signal processor 14 comprises a servo control section 14a and an information signal reproducing section 14b. The servo control section 14a is coupled to the disc drive motor 15 and the pickup device 17 for supplying servo control signals thereto, as described later in brief. The information signal reproducing section 14b is coupled to the pickup device 17 for receiving an RF signal (RF; abbreviation of Radio Frequency) Srf read from a CD 18 and sent to the CD reproducing mechanism section 14.\nThe disc drive motor 15 drives the disc turntable 16 under the servo control of the servo control section 14a. The pickup device 17 shifts in the radial direction of the CD 18 under the servo control of the servo control section 14a. Thus, the pickup device 17 traces a record track on the CD 18 in a CLV (Constant Linear Velocity) speed. The pickup device 17 reads the CD 18 during the tracing of the record track of the CD 18 so that an RF signal Srf corresponding to the digital data recorded on the CD 18 is output from the pickup device 17.\nThe pickup device 17 outputs the RF signal Srf corresponding to the digital data, as described above. The RF signal Srf is applied to the information signal reproducing section 14b of the signal processor 14. A data slice circuit (not shown) in the information signal reproducing section 14b slices the RF signal Srf at a prescribed threshold level so that the RF signal Srf is shaped to the original digital data.\nThe servo control section 14a of the signal processor 14 generates two servo control data SCm and SCp and a bit synchronizing signal BS from the reproduced digital data. A motor servo control data SCm is applied to the disc drive motor 15, as described before. A pickup servo control data SCp is applied to the pickup device 17, as described before. The servo controls for both the disc drive motor 15 and the pickup device 17 are carried out in using the digital data obtained from the output of the pickup device 17, i.e., the RF signal Srf, and also the clock pulse Pck applied from the clock generator 13.\nThe information signal reproducing section 14b of the signal processor 14 demodulates an analog data AD from the reproduced digital data under the control of the clock pulse Pck. The analog data AD is applied to the cassette tape recorder 11. The cassette tape recorder 11 then records the analog data AD.\nThe conventional reproduction apparatus of digital discs, however, have a drawback that a reproduction of the digital discs takes a relatively long time the same as a standard record time provided for each of the digital discs."}
-{"text": "The present invention relates to apparatus, including a cartridge adapted to be removably mounted to the apparatus, for severing a tubular workpiece.\nTube cutting apparatus employed to sever sections of a continuous length of tubing or similar elongated stock fed from a source thereof such as a mill is known. Tube cutting apparatus in which a cut-off die set, driven by a ram assembly and clamped to the continuous length of tubing prior to and during the severing operation is exemplified by U.S. Pat. 4,653,368 issuing Mar. 31, 1987 to Riera et al., the disclosure of which being incorporated by reference herein. Riera discloses a tube cut-off die set including mechanism for providing a horizontal notching cut in a tubing wall followed by a vertical severing cut as the tube is fed through the set. The tub cut-off die set includes a movable upper shoe, a stationary lower shoe, guide posts and associated bearings affixed to the shoes and a ram providing the relative reciprocation between shoes. A tube holding arrangement therein comprises two axially spaced pairs of complementary die jaws being mounted on the lower die show for releasably clamping the tube. Downward reciprocation of a cam driver extending from the upper die shoe engages a pair of cam rollers associated with the lower die shoe to drive the die jaws about the tube and the blades between the jaw members.\nThe die set is used to cut sections of tubing as the tube is continuously manufactured. While suitable for its intended operation, the above-referenced die set does not allow for rapid adjustment of the die jaws so as to change their grip about the tube, and does not permit rapid replacement of the die jaws to accommodate different tube sizes and shapes. Since each tube diameter and cross section requires different jaw members it would be desirable to have jaw members that can be replaced and preset at a location spaced from the die set so that there would be no downtime of the tube-cutting apparatus. Further, a more compact jaw mounting arrangement would be desirable to reduce the overall size of the die set.\nIt is an object of this invention to provide a tube cut-off machine that maximizes the ease with which die jaws are replaced, to accommodate different tube diameters, without stopping the machine operation, and provide the user with an adjustment arrangement whereby the clamping force by the die jaws on the tube can be easily adjusted. Advantageously the use of a self-contained cartridge that includes the jaw members, cam followers and adjustment arrangement would allow the user to keep the machine cutting one particular diameter size while other cartridges are retrofitted with the appropriate jaw members, thereby reducing downtime of expensive machinery when it is time to change the jaws.\nPart of the setup time of jaw members associated with the Riera machine was removing the cam driver prior to removal of jaw holders and then remounting the jaw holders and cam driver. This arrangement has an element of machine downtime. It would be desirable to provide rapid connecting and disconnecting arrangement for the cam driver. Desirably, such an arrangement would associate the desired jaw member spacing with the cam driver.\nA further object of this invention is to provide a replaceable kit for use in a tube cutting machine which is simple, has jaw holders and jaw members preset, and is insertable as a unit. Advantageously, because the costs associated with the downtime of a machine is high and the costs of a cam driver relatively small, teardown and setup costs can be more than justified by such a kit. In this regard the kit advantageously permits only minimal downtime between removing cartridges and associated jaw holders when different tube diameters are to be cut.\nOne problem sometimes experienced in connection with tube severing machine is that the severed tube end portion is locally deformed and has what are termed \"dimples.\" Tubes that are so formed are waste and thus costly. While the Riera construction including a horizontal notching blade cantilevered at the end of a cross slide in connection with the vertical severing blade has worked well in most situations, tube dimpling with some cross sections and tube thicknesses results. A tube severing arrangement which does not produce deformed cross sections would be desirable.\nAnother object of this invention is provision of a tube severing arrangement that permits severing of thicker and larger diametered tubes without breaking the blades or increasing the size of the machine into which they are mounted. Further, configuring the notching blade so as to have an offset portion surprisingly allows such tube cross sections to be cut.\nThe foregoing objects and advantages are accomplished by a tube cut-off machine for severing elongated material moving continuously longitudinally of its length, comprising relatively reciprocable upper and lower die shoes, one shoe mounting a tube cut-off blade and a cam driver, and tube clamping means for releasably clamping the tube during the severing operation. The tube clamping means includes two laterally spaced jaw holders each carrying a jaw member and positioning the tube clamping surfaces of the jaw members in confronting relation and a cam follower operably connected to one jaw holder.\nIn accordance with this invention, the tube clamping mans is characterized by a compact, self-contained, cartridge that is removably mounted to the die set as a single by replaceable unit. The cartridge comprises a carriage adapted to be mounted on the mounting surface of the lower die shoe so as to be stationary relative thereto, a pair of laterally spaced cam followers each being mounted on a cam holder one and the other follower being a roller and a replaceable wear pad, the two jaw holders and cam follower holders, respectively, being laterally spaced and slidably disposed on the carriage in confronting relation, bias means for laterally biasing the jaw holders apart, and adjustment means including a frame having an adjustment screw and operably associated with the carriage for adjusting the relative lateral separation between the jaw members, the frame maintaining the holders in side-by-side relation and the screw adjusting wear in the wear pad.\nAdvantageously the holders are separately manufactured and thus the assembly is less expensive because the need for a \"beefed up\" carriage construction is minimized. Separate parts weigh less, thus reducing forces needed during clamping/unclamping of the tube to overcome friction and inertia. Advantageously, the wear plate absorbs shock loads which might otherwise shear the pivot pin supporting a cam roller for rotation, resulting in costly repair and machine downtime.\nA kit adapted to be mounted on the die set comprises the cam driver being operably disposed within the cartridge between the cam followers and their spacing preadjusted so that the jaw members are at their desired separation for clamping about the tube. The kit is mounted separately to the machine and requires minimal assembly time.\nIn accordance with one mounting arrangement, the cartridge is slid beneath a pair of gibs to secure the cartridge to the lower shoe, a dowel pin in a driver holder is connected to the upper shoe, and a pair of driver clamping plates on the driver holder are tightened behind the cam driver.\nThe cam driver has compliant load bars adapted to engage both cam followers to reduce uneven loads placed on the cam followers.\nFurther in accordance with this invention, the die set includes four bushings with their vertical axes disposed in a rectangular array with first, second and third adjacent pairs being secured together, respectively, by a first tie bar, a second tie bar and rail mounting plate, and the horizontal notching blade is vertically adjustable relative to a holding bar adapted to move between vertical planes including the first and third pair of bushings. Advantageously, a stiffening arrangement including three guide bushings including the second tie bar and rail mounting plate, the bashing axes forming a right triangular array, can stiffen the die set to obviate wobble of the notching blade."}
-{"text": "Biodegradable softeners have attracted recent attention in the prior art. For instance, in German patent no. 197 43 687, in the name of Henkel KGaA, readily biodegradable detergents are described, which contain oligomeric esterquats obtained by quaternizing oligoesters of mono and di-carboxylic acids in combination with alkylene oxide adducts on fatty acid amines.\nThe international patent application WO-A-01/47489, in the name of Cognis Deutschland GmbH and Bigorra Llosas et al., discloses fiber brightening and softening agents comprising esterquats obtained by reacting alkanol amines with mixtures of fatty acids and dicarboxylic acids, optionally alkoxylating the resulting esters, and quaternizing the products; and auxiliary materials selected from (non)quaternized fatty acid amides, betaines, nonionic surfactants, polyols and/or their derivatives, alcohols and/or hydrotropes.\nIn the European patent application 1 136 471, in the name of KAO Corporation S.A., alkanol amine esters are described which are based on the esterification reaction of alkanolamines, carboxylic acids and fatty alcohols. The alkanolamines and fatty alcohols are optionally alkoxylated. In addition, the cationic surfactants and esterquats obtainable therefrom are disclosed.\nThe cationic surfactants and the esterquats disclosed in said European patent application have a high degree of biodegradability, but compared to the biodegradable esterquats of the prior art also exhibit a high degree of efficacy in softening and conditioning natural and synthetic fibers, such as hair, or fibers used in textiles and paper.\nIn a further aspect, said European patent application relates to aqueous fabric-softening compositions which contain the cationic surfactants or esterquats, optionally together with other active softening substances. Particularly, these fabric softening compositions contain, in an aqueous medium optionally containing constituents selected from those normally used in fabric softener composition: (a) cationic surfactants or esterquats obtainable from the alkanolamines described, (b) one or more cationic surfactants which are active as fabric softeners, and (c) one or more non-ionic fabric-conditioning surfactants, wherein the amount of (a)+(b)+(c) is 2-60 wt. % based on the total composition; the amount of (a), based on the total of (a)+(b)+(c), being 2-100 wt. %; the amount of (b), based on the total of (a)+(b)+(c), being 0-98 wt. %; and the amount of (c), based on the total of (a)+(b)+(c), being 0-40 wt. %."}
-{"text": "This invention relates to an improved transducer for use in a vibratory viscometer; which transducer is particularly useful for, but not limited to, the inline detection under flow conditions of viscous and elastic properties of fluids being processed.\nVibratory viscometers are well known in the art and generally employ a transducer which has an immersible portion which is vibrated with a small amplitude. Fluid viscosity/density/viscoelasticity can be determined from the frequency and/or amplitude changes in the vibration and/or the power required to sustain the vibration when the immersible portion of the transducer is immersed in a fluid.\nSuch transducers generally comprise (i) an immersible tip, (ii) an electromagnetic drive and (iii) an electromagnetic or piezoelectric pickup. Transducers of this type are described by J. G. Woodward in \"Vibrating Plate Viscometers\", Electronics, February 1952; and by J. D. Ferry in \"Viscoelastic Properties of Polymers\", published by John Wiley & Sons, New York, 1970.\nPrior Oscillatory Viscometers\nThe oscillatory viscometers described below exemplify the prior art, and have transducers which interact with the fluid being measured by \"surface loading\" of an oscillating portion of the transducer; said term having been used J. D. Ferry in his aforementioned publication.\nThe oscillating surface generates shear waves in the liquid or other fluid. Viscoelasticity is measured in terms of the characteristic impedance of the liquid.\nIn the vibrating plate viscometer described in the aforementioned publication of J. G. Woodward, the moving end of a vibrating reed supports a plate which detects viscous resistance of the fluid in which it is immersed. The reed is driven electromagnetically and the oscillations are picked up by a barium titanate piezoelectric block. The reed is clamped at one end, exposing the driver and pickup to fluid vapors. The viscometer measurements are confined to viscous loss determined from the decrease in amplitude of oscillation observed when the plate is immersed. The plate vibrates in a direction essentially perpendicular to its major surfaces.\nA. Konno, S. Malsino and M. Kameko, in Japan Journal of Applied Physics 189 (1968), reported their measurements of storage modulus and viscous loss by oscillating an immersed very thin microscope slide at 100 Hz. This apparatus, driven by a moving coil, is reported as strictly a research instrument. The internal workings of the transducer were exposed to fluid vapors.\nA viscometer known as Le Viscosimetre \"Pivert\" is sold by the Societe Francis de Service, 8 rue Nobel Zl, 45700 Villemander, France. The sensor tip of this viscometer (also called \"Sofraser\") is a U-shaped stainless steel needle. One leg of the U-shaped needle is mechanically driven sinusoidally at 125 Hz., which is near the mechanical resonance frequency of the transducer. Both legs of the needle are clamped at nodal points. Frequencies are measured at both legs. The phase difference between the frequencies is transformed into a voltage or current related to the viscosity of the liquid in contact with the needle. The manufacturer states that installation conditions can decrease accuracy and therefore should be carefully checked to be sure of optimum accuracy. The range of viscosity measurement is 0.5 to 15,000 mPa.s (cPs). The environmental limits of utilization are 200.degree. C. and pressure to 100 bars.\nThe Dynatrol Viscosity Detector is sold by Automation Products, Inc., 3030 Max Roy, Houston, Tex. 77008. Although this instrument does not employ a blade sensor, it does employ welding at the node point of a resonating system in order to isolate the oscillating probe from the driver and pickup. Viscosity-density product is detected by an immersed 5 inch long stainless steel probe. The probe is made from a stainless steel rod and bent very much like a hairpin. Both ends of the rod are welded at nodal points to a supporting plug. One end of the probe penetrates into the housing where it is electromagnetically driven to resonate in flexure at 120 Hz. The other end of the probe also penetrates through the plug into the housing where there is an electromagnetic pickup coil. The decrease in amplitude of vibration of the immersed probe due to interaction with the fluid is electronically converted to viscosity-density product.\nThe Labor-Viskosimeter QV35 is marketed by Bopp & Reuther GmbH, Car-Reuther-Strasse 1, Postfach 310140, D-6800 Mannheim 31, West Germany. This instrument uses an oscillating transducer to measure viscosities of laboratory samples of liquids by means of a quartz crystal sensor oscillating in torsion at 55 KHz. The damping effect of the fluid being measured on the amplitude of oscillation is converted into viscosity. The temperature range of this instrument is limited to -50.degree. to 150.degree. C.\nThe Model 1800 Viscometer sold by Combustion Engineering, Inc., P.O. Box 831, Lewisburg, W.Va. 24901 utilizes an oscillating sensor blade. Short pulses of current are applied to the top of the blade, which is composed of a magnetostrictive alloy. The blade protrudes through a metal diaphragm. Each pulse causes the blade to vibrate at its natural frequency of 28 KHz. When the amplitude of vibration of the immersed blade has fallen to a preset value that relates to the viscosity of the liquid, another pulse is automatically applied. The change in pulse rate is proportional to the square root of viscosity-density product. The range of viscosities measured is from 0 to 5,000 centipoise x grams/cm..sup.2 in ranges of 0-50, 0-500, and 0-5000. This viscometer was produced for many years by Bendix Corporation, and was first described as the Ultra-Viscoson by W. Roth and S. R. Rich in Jr. Applied Phy. 24 940-950, July 1953. The Bendix Ultra-Viscoson is described on page 308 of Viscosity and Flow Measurement by S. R. VanWaser, S. W. Lyons, K. Y. Kim and R. E. Cowell, Interscience Publishers, New York, 1963. Shortcomings included very high frequency of measurement, fragility of very thin strip (blade), lack of sensitivity at very low viscosities, need to flick the strap from time to time to relieve strains in the magnetostrictive alloy; and since strips were easily bent, replacement blades needed to be available.\nHermetic sealing between driver and pickup is a feature of the \"Vibrating Sphere\" and \"Viscoliner\" oscillating viscometers of Nametre Company, 101 Forrest St., Metuchen, N.J. 08840, the assignee of the present application. However, the spherical and cylindrical sensors are somewhat obstructive to flowing fluids, particularly slurries. In order to reduce the obstruction to flow, the diameter of the pipe in which the fluid flows needs to be increased so as to satisfactorily accept the oscillating sensor. Since the mode of oscillation has two degrees of freedom, care must be taken to chose between in-phase and out-of-phase torsional motion. Further, the oscillating surface generates diverging shear waves that may be so long in wavelength that they are not conveniently reflected for accurate measurement of highly viscoelastic fluids.\nU.S. Pat. No. 4,729,237 describes a tuning fork transducer, each of the two arms of the fork having welded to it a blade that is oscillated in the liquid. It is claimed that one blade on one arm of the tuning fork gave less accurate viscosity measurements than having blades on each arm. This viscometer is a laboratory instrument. The specification describes vertical motion of a sample container to immerse the two blades in the liquid.\nObjects of the Present Invention\nFor a long time there has been great need for small inline blade sensors that do not impede flow and that are intrinsically separated from the fluid being measured. Examples of these needs include transport of mineral slurries such as powdered coal and powdered lime where particle concentration must be controlled. The consistency of food fluids such as ice cream, coffee and bread dough needs to be monitored and controlled by means of reliable rugged viscosity sensors. Better control of viscoelastic properties of polymer fluids can be achieved by employing suitable sensors. Accordingly, an object of the present invention is to provide an inline blade sensor for a vibratory viscometer that impedes fluid flow to a substantially lesser extent than prior art sensors. Another object of the invention is to provide such a sensor in a configuration that facilitates isolation of the drive and pickup portions of the transducer from the fluid.\nOther objects of the invention are to provide transducer-sensor devices that are relatively small, are rugged, easily inserted into industrial process pipes and equipment; that operate over a broad range of viscosities under hot and cold conditions and at high pressures; that can be speedily installed and quickly operated in remote and dangerous locations; that are dependable over long stretches of time; that can be easily removed, cleaned and reinstalled; that are capable of providing viscous loss and/or transducer sensor frequency signals for conversion into loss modulus and storage modulus values; that are capable of being vibrated over ranges of frequencies as well as at mechanical resonance; and that are adaptable for use in very viscous and very elastic liquids."}
-{"text": "This invention relates to a remote control apparatus and, more particularly, to an axial alignment aid and related method for facilitating remote control operations.\nMany attempts have been made to remotely align two objects in space. For example, the National Aeronautics and Space Administration (NASA) has done extensive work in the area of docking spacecraft. As discussed in \"Communications, Tracking and Docking on the Space Station,\" by Erwin et al., published by the Institute of Electrical and Electronics Engineers (1982), NASA has proposed the use of a triangular array of reflectors on a passive target to define a plane and facilitate docking of an active vehicle thereto. Sequential tracking of each reflector from the active vehicle in space allows altitude determination between the active vehicle and the passive target. However, this triangular reflector system necessarily employs computers, thus making it extremely complex. In addition, the system is incapable of accurately and remotely aligning an axis of importance of the passive target to an axis of importance of the active vehicle to allow \"soft\" docking.\nOther attempts have been made to accurately align or dock one object relative to another. For example, U.S. Pat. No. 3,269,254, issued to COPPER et al., discloses an optical apparatus including: a light source; a beam splitting cube; a lens system and a grated mirror, all on one side of the light source; and a measuring reticle having coordinates thereon located on the opposite side of the light source. Light from the light source is directed by the beam splitting cube through the lens system. The light then reflects off of the mirror to again pass through the lens system and the beam splitting cube. Finally, an image is formed at the measuring reticle to provide information regarding orientation of a body o which the mirror is attached, as defined by azimuth, pitch and roll. Because the mirror described in this patent is flat, if the beam of light directed at the mirror is too far off-axis, the light beam can be reflected too far off, thus preventing the formation of an image at the measuring reticle. As a result, an accurate off-axis measure is not available with this patented device. Further, this patented device is sensitive to distance, i.e., the effectiveness of the device is lessened as the measuring reticle and the mirror are placed farther apart. Finally, this device is not adapted to allow docking of one member relative to another member.\nIn addition, U.S. Pat. No. 2,352,179, issued to BOSLEY, uses a set of four photo-electric cells as detectors to provide information as to the orientation of a lens. That is, the set of cells indicates from one perspective how far away one is from a target; but such detectors cannot accurately indicate the axis that is defined by the target. Therefore, such a system using a set of four detectors is also incapable of allowing on-axis docking alignment of one object relative to another.\nOne application where remote control docking would be particularly important is in nuclear powered steam generators. For example, performing a task within the steam generator with a tool or end effector and a robot for receiving the end effector is currently performed as follows. A first set of cameras positioned in the steam generator allow a conventional robot arm known as ROSA (remote operation service arm) to be remotely viewed and moved to the general vicinity of one or several end effectors positioned in the steam generator. These cameras, however, cannot accurately and automatically align the respective axes of the chosen end effector and ROSA, i.e., cannot reliably effect docking and attachment.\nThe end effector also has a camera mounted thereon which ultimately is used to remotely view the worksite during performance of the task via a conventional feedback system using a remote viewing station, robotic controls and means for moving ROSA.\nAs a result, a man must be introduced into the steam generator to manually align and attach the end effector to ROSA. However, due to the potential radiation exposure, tedious alignment and attachment steps necessary, and the relatively heavy members being manipulated, it is desired to make the alignment and attachment of end effectors to ROSA in the steam generator an entirely remotely controlled operation having high reliability A more general need also exists for a means by which accurate alignment can be made between two important axes of respective docking members using single perspective visual guidance."}
-{"text": "(a) Technical Field\nThe present invention relates to a control method for improving fuel efficiency of a hybrid electric vehicle. More particularly, the present invention relates to a method for controlling a fuel cut operation to be performed in a wider transmission region by a reaction control of a motor, thus improving fuel efficiency of a hybrid electric vehicle.\n(b) Background Art\nIntensive research has been underway to provide a high efficiency hybrid electric vehicle (HEV) by improving fuel efficiency and emission performance.\nOne of the various techniques for improving fuel efficiency and emission performance is to provide an HEV with idle stop system, in which the engine of the vehicle is stopped while the engine is idling as the vehicle is stopped.\nAnother technique for improving fuel efficiency and emission performance is to increase fuel cut operation. A fuel cut mode saves energy by stopping fuel injection when a vehicle running at a high speed is decelerated.\nHowever, it has a problem of the fuel cut range being restricted. More specifically, a sudden change in engine torque can occur when the engine is restarted as the fuel is injected. When a torque is applied at a low speed, for example, at a low gear 1 or 2, the vehicle may be shaken due to the sudden change in engine torque.\nFor example, when the engine speed is lowered below about 1500 RPM, a fuel in (fuel re-injection) is performed. Although it is desirable to increase the fuel cut period for improvement of fuel efficiency, impact is brought about at the moment when the fuel cut operation is cancelled (i.e., rewetting is started). To improve fuel efficiency, it is necessary to prolong the fuel cut period at a lower speed and RPM.\nIn the event that a fuel cut operation is performed for an ordinary vehicle, the fuel cut operation is cancelled at about 1500 RPM to revive the engine. Converting the RPM into a vehicle speed, the fuel cut operation is cancelled at about 60 km/h (and 1500 RPM) in case of a vehicle not exceeding 2000 cc, although it may vary according to a gear ratio of the vehicle.\nAccordingly, if the vehicle is running at 60 km/h or higher for a predetermined period of time, the fuel cut effecting condition is permitted. The fuel cut operation continues during deceleration period and is cancelled when a preset condition is met.\nIn detail, the fuel cut cancellation is made when it reaches a predetermined RPM with respect to coolant temperature, and a corrected RPM is considered thereto. Moreover, in the case where an air conditioner is turned on, an additional RPM is considered and, in the case where the air conditioner is turned off, the fuel cut cancellation is made at about 60 km/h and 1500 RPM.\nHEVs developed so far reach about 1500 RPM at 60 km/h and the vehicle speed is decreased while maintaining about 1200 RPM up to about 15 km/h. With the current fuel cut logic for gasoline vehicles, it is thus impossible to improve fuel efficiency of HEVs. Moreover, a torque of about 15 to 25 Nm is suddenly applied during the fuel cut cancellation, which results in a serious impact that a driver can feel when the gear ratio is low.\nThe information disclosed in this Background section is only for enhancement of understanding of the background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art that is already known to a person skilled in the art."}
-{"text": "This invention in general relates to the void fraction measurement in gas/liquid two-phase flow system. Particularly, the invention relates to a method and a pressure-differential device to measure the void fraction of a gas/liquid two-phase stratified flow within a horizontal pipeline system. The said method and said device are applicable to flow at steady and transient states, and at high-temperature/pressure conditions.\nVoid fraction is a crucial parameter in the determination of energy and momentum transfer in a two-phase flow system. In usual practice, the void fraction is defined two-dimensionally, in reference to a chosen cross section of the pipe under measurement, as the area of gas phase passing through that cross section divided by the whole cross-sectional area. At present, the instruments most often used for void fraction measurement are the expensive and complicated xcex3-ray and \"khgr\"-ray densitometers, in which void fraction is derived from radiation attenuation of one or several xcex3-ray or \"khgr\"-ray beams traversing through the gas/liquid two-phase flow. FIG. 6 illustrates such a densitometer commonly used in industry, which includes radiation source 60, collimators 61, two-phase flow measurement section 62, radiation shields 63, and a detector and counting system 64. This layout was first reported by Edelman et. al. in International Journal of Heat and Mass Transfer, vol. 28, no. 7, 1985, where xcex3-ray is applied to measure the average cross-sectional void fraction of a steady-state boiling flow within a stainless steel pipe. FIG. 7 is another example, presented by Kagawa et. al. in International Topical Meeting on Nuclear Reactor Thermal Hydraulics, NUREG/CP-0014, vol. 2, 1990, where \"khgr\"-ray scanning method is applied to a high-temperature/pressure flow system to measure the distribution of cross-sectional void fractions in a horizontal pipe and the system is at transient state due to a severe leakage somewhere in the pipeline system. This is an even more complicated instrument than the xcex3-ray densitometer just mentioned, in that, more than those basic component parts shown in FIG. 6 the instrument has, referring to FIG. 7, a radiation scanning location-control mechanism 65 and a radiation scanning beam-control electronic system 66. The measurement accuracy of xcex3-ray or \"khgr\"-ray densitometer depends on the degree of homogeneity of the gas/liquid two-phase flow: better homogeneity of the fluid obviously promises more reliable measurement result. However, it is well known that, the void fraction of a gas/liquid two-phase flow seldom is, if ever, homogeneously distributed. The measurement accuracy of a radiation densitometer also depends on the incident angle and location of each radiation beam, because each beam, different from the rest, is specifically attenuated by a certain wall thickness of pipe and by a certain portion of gas/liquid two-phase flow. As a rule, the farther away is a radiation beam from pipe axis, the less reliable is the measurement data. Moreover, depending on the material and thickness of pipe wall, the radiation attenuation may be so great as not to leave any useful information for the detector to resolve. For a specific example, one may consider the case of a horizontal natural circulation of a gas/liquid two-phase thermal-hydraulics system, in which the density difference between low flow rate liquid phase (e.g. water) and high flow rate gas phase (e.g. vapor) is usually so large that, under the influence of gravity, the two phases virtually separate and become two quite distinct flows, namely the upper layer gas flow and the lower layer liquid flow, thus the term flow stratification. Neither xcex3-ray nor \"khgr\"-ray densitometer is at all suitable for such flow condition for the want of fluid homogeneity.\nTo overcome the aforesaid deficiencies of radiation method for void fraction measurement, the present invention, taking advantage of the fact that flow stratification phenomenon often occurs in natural circulation of gas/liquid two-phase flow within horizontal pipeline system, discloses a method to determine the void fraction of gas/liquid stratified flow; wherein a couple of pressure-signal tubes are vertically flush-mounted to the top and bottom interior surface of the horizontal pipeline to measure the pressure differential across the pipe. A thermocouple-rake temperature probe assembly is used to obtain the average temperature of the liquid phase of a gas/liquid two-phase stratified flow, in which the thermal stratification also occurs. Therefore, the average mass density of the liquid phase of the stratified flow is obtained. Then the liquid-phase level of the stratified flow is calculated from the aforesaid pressure-differential data and liquid-phase mass density data. Finally the void fraction is derived from the known geometric relation between said liquid-phase level and the cross section of the pipe.\nSince commercialized instruments for void fraction measurement of gas/liquid two-phase flow, such as xcex3-ray or \"khgr\"-ray densitometer, are far from satisfactory when flow stratification occurs, owing to the reasons briefly discussed in the foregoing, it is the principal objective of this disclosure to resolve this problem by providing a novel method for void fraction measurement, applicable to vapor/water, or any other kind of gas/liquid two-phase stratified flow within a horizontal pipeline system. According to this disclosure, a couple of pressure-signal tubes are vertically flush-mounted to the top and bottom interior surface of the pipe at a selected location, so as to measure the pressure differential across the cross section at that location; and, through liquid-phase mass density correction, this pressure-drop data is transformed into the liquid-phase level within the pipe; and the said liquid-phase level, by its geometric relation with the cross section of the pipe, is then transformed into void fraction. The device herein disclosed can accurately determine the void fraction of a horizontal gas/liquid two-phase stratified flow, be the flow at steady or transient state, and whether or not the flow is at high-temperature/pressure condition.\nWhen the flow being measured is at high-temperature/pressure condition or is at transient state, physical effects like gravity and heat that may render pressure-differential measurement difficult must be put into consideration. Not only does gravity produce flow stratification, but it may also pull down some pressure-signal transmitting medium from the upper pressure hole, resulting in voids or even void sections present within the upper pressure-signal tube. To obviate the gravitational effect four measures are taken: (1) keep as small as possible the diameter of pressure-signal tube that is near the pressure hole, (2) keep as short as possible the vertical section of the upper pressure-signal tube that is near the pressure hole; (3) keep as small as possible the two locations at top and bottom pressure holes where pressure-signal transmitting medium is present; and (4) have pressure-signal tubes flush-mounted to the top and bottom interior surface of the pipe. As to heat, obviously it can increase the temperature of pressure-signal tubes, especially at the sections near the pressure holes, because of thermal conduction. A density gradient will therefore occur within pressure-signal transmitting medium. This problem will be especially serious at the upper vertical section of the pressure-signal tube, and will greatly affect measurement accuracy of the liquid-phase level of a gas/liquid two-phase stratified flow within a horizontal pipeline system. If accident such as leakage or cleavage occurs in a high-temperature/pressure closed system, resulting in pressure decrease faster than temperature decrease, vaporization of the liquid phase within the system may begin rapidly (this is called flashing phenomenon), and this may sometimes lead to vaporization of pressure-signal transmitting medium near the pressure holes. Flashing of pressure-signal transmitting medium may also result in voids or even void sections within pressure-signal tube. Once voids or void sections are present within the vertical section of pressure-signal tubes, any further measurement is not merely futile but, in some cases, possibly dangerous may result. This is because the measurement device is no longer reliable: it may produce false signals which may accordingly trigger the safety systems of, say, a nuclear reactor to work, or, contrariwise, it may ignore accident evolution and delay the safety systems, and neither case is exactly desirable. To mitigate the effects of thermal conduction and flashing phenomenon, this invention has each one of the pressure holes and its nearby space encased within a cylindrical coolant cladding to keep the temperature of the pressure-signal transmitting medium at safe condition. Experiments done in this laboratory prove that this local cooling strategy can effectively prevent loss of pressure-signal transmitting medium due to density gradient or vaporization. Thus, more than providing a method and a device to measure the void fraction of a gas/liquid two-phase stratified flow within a horizontal pipeline system, this invention provides a preventive method and a device to effectively overcome the problem of loss of measurement accuracy due to loss of pressure-signal transmitting medium encountered in transient flow system by the commonly used pressure-based instruments such as fluid-level meter, flowmeter, pressure-drop meter, pressure meter, etc."}
-{"text": "In current digital wireless systems, the traditional up-conversion chain (or significant portion thereof) is primarily analog and includes types such as super-heterodyne, low intermediate frequency (IF) and zero IF up-conversion technology. These technologies start with the conversion of inherently digital signals to analog signals through high performance digital-to-analog (D/A) converters, generally due to the higher frequencies involved. Once converted to the analog domain, various combinations of analog filters, amplifiers, mixers and modulators (and perhaps other analog elements) are cascaded to achieve the up-conversion from the output of the A/D converter(s) to the radio frequency (RF) band of interest (transmit RF signal).\nLikewise, on the receiver side, the traditional down-conversion chain (or significant portion thereof) is primarily analog including such types as super-heterodyne, low IF and zero IF down-conversion technology. To achieve the down-conversion, various combinations of analog filters, amplifiers, mixers and demodulators (and perhaps other analog elements) are utilized to achieve the conversion from the RF band of interest (receive RF signal) to the input to A/D converter(s).\nComponent variation, tolerances, and aging all affect the design requirements, costs, and manufacturability of the analog up-conversion (transmitter) and down-conversion (receiver) chains. Accordingly, there is needed a digital transmitter and digital receiver that utilizes digital technology for the up-conversion and down-conversion chains."}
-{"text": "This invention relates to computer systems and circuits; and more particularly, to methods and apparatus for generating an appropriate reference voltage level for differential signaling.\nComputer systems are consistently being updated to operate at faster rates. In doing so, many of the components of these system must operate within new parameters, or be replaced at an additional cost. Achieving these higher operating rates further requires devices within these systems to communicate among themselves at faster rates.\nA technique employed to increase the rate at which two devices can communicate with each other is decreasing the voltage swing between a communication signal\"\"s high and low values. Typically, this includes reducing the supply voltage level at which these devices operate. This decrease in voltage swing allows a signal to transition faster between values. However, as this voltage swing is decreased, additional precautions must be taken to ensure that noise does not interfere with the signal.\nDifferential signaling is one method used in communicating between devices in a system and is useful in avoiding errors induced by noise. Typically, each sending device provides a voltage reference and a communication signal which are typically connected to a differential amplifier. Because approximately the same noise is typically induced on both the voltage reference and a communication signal, a receiving differential amplifier produces a relatively xe2x80x9ccleanxe2x80x9d communications signal.\nHowever, not all sending devices can provide an appropriate reference voltage. For example, when a system is upgraded or a new version is introduced that requires a different reference voltage level and possibly uses a different supply voltage, an older component of the system might not be able to provide the new reference voltage.\nSuch is the problem introduced by the Accelerated Graphics Port (xe2x80x9cAGPxe2x80x9d) Interface Specification Revision 2.0, May 1998, which requires a universal AGP target (such as the 82465GX GXB manufactured by Intel Corporation) to support both 3.3 Volt and 1.5 Volt operation on the AGP interface supply voltage (Vddq). The specification stipulates a voltage reference that is nominally 0.5*Vddq for 1.5 Volt operation, and nominally 0.4*Vddq for 3.3 Volt operation. Therefore, the value of the reference voltage must be adjusted to match the requirements based on the supply voltage used. What is needed is a system for efficiently providing the appropriate reference voltage.\nAccording to the invention, systems, apparatus and methods are disclosed for providing one or more bimodal reference voltages for using in differential signaling. Typically, a system comprises a first component requiring a first reference voltage; a second component requiring a second reference voltage; a power source to produce a supply voltage; a first voltage reducer electrically coupled to the power source, the first voltage reducer to produce a third reference voltage using the supply voltage; a multiplexor electrically coupled to the first voltage reducer to produce the first reference voltage by selecting between at least the third reference voltage and a fourth reference voltage; and a second voltage reducer electrically coupled to the power source and to the second component, the second voltage reducer to produce the second reference voltage using the supply voltage."}
-{"text": "1. Field of the Invention\nThe disclosed concept relates to a switchgear and, more particularly, to a switchgear wherein an enclosure housing assembly includes an extendable conductor assembly and a ground bus bar assembly that does not extend over the path of travel of a movable carriage assembly.\n2. Background Information\nElectrical switching apparatus, such as but not limited to circuit breakers, are often disposed in an enclosure or housing assembly. The electrical switching apparatus is disposed on, or is incorporated with, a movable carriage assembly. When the electrical switching apparatus need repair or other maintenance, the electrical switching apparatus is drawn, at least partially, from the enclosure. The removal of the electrical switching apparatus from the housing assembly may, or may not, decouple the electrical switching apparatus from the power source (hereinafter \u201cline\u201d) and/or the load. For the safety of the operators, a movable carriage assembly ground connection, i.e. a grounded conductor, is required by various regulations and safety procedures. For example, the frame of the circuit breaker is, typically, an electrically conductive material which needs to be coupled to ground whenever secondary wiring connections mate and/or become energized with electrical power\nThe ground connection allows the electrical switching apparatus to be in electrical communication with the ground continuously from the fully inserted/connected position until the electrical switching apparatus is drawn out a safe distance from the housing assembly. The purpose of this grounding feature is to protect operators from injury as a result of an electrical issue in the secondary control voltage system, which is typically less than 250 volts. The ground connection is also used with Ground and Test Devices that connect the primary voltage terminals, which typically can carry up to 15,000 volts, to ground so that the operators are protected while working on the conductor buses.\nTo allow for such a continuous ground connection, the housing assembly includes a large conductive bar (hereinafter \u201cground bar\u201d), such as, but not limited to a copper bar, secured to the housing assembly. Further, the electrical switching apparatus includes a sliding electrical coupling that is structured to be slidably coupled to, and in electrical communication with, the ground bar. That is, the ground bar extends over the path of travel of the electrical switching apparatus so that as the electrical switching apparatus moves out of, or into, the housing assembly, the sliding electrical coupling moves along the ground bar maintaining the electrical switching apparatus in a grounded state.\nBecause advancing technology continually reduces the overall size of an electrical switching apparatus and the associated housing assembly, physically locating large conductive bars that extend the full length of the housing assembly is becoming increasingly difficult. Additionally, market demand for copper makes large ground bars expensive. By employing a short ground bar at the rear of the housing assembly, testing equipment can still effectively connect to, and provide the shunted path for, primary voltage and current; however, the requirement for a continuous connection throughout the draw-out range still must be met.\nThere is, therefore, a need for a housing assembly that provides for a grounded connection over the length of the draw-out path of travel but which does not include a long ground bar. There is a further need for such a housing assembly to accommodate existing electrical switching apparatus."}
-{"text": "In computer graphics, traditionally colors are represented using a combination of primary colors. For example, in the RGB color space, the colors Red, Green and Blue are blended together to get a range of colors. In order to manipulate images in the RGB color space, scripting languages are designed to allow the users to create scripts that describe image-processing operations at a high level. For example, a script can be written to handle transition between two video clips or between one video clip and one still video. The scripts do not necessarily work only with transitions. They may also be written to work with single source clips. For example, a script that takes an existing clip of video and drop colors on it. The scripts may also be \u201cgenerators\u201d which have no inputs but create an output.\nFinal Cut Pro (FCP) is a movie editing and creating software produced by Apple Computer, Inc. of Cupertino, Calif. (\u201cApple\u201d). It includes a scripting engine. The FCP scripting engine allows users to write scripts which perform various image manipulations, ranging from basic operations such as \u201cblend\u201d \u201cchannelfill\u201d and \u201cmultiplyChannels\u201d to more complex operations such as \u201clevelmap\u201d and \u201ccolorkey.\u201d Of course, most scripts combine more than one of these operations to build interesting effects or transitions.\nThese functions operate on one or more images. Consider these examples. \u201cChannelfill\u201d fills in one or more color channels of a destination buffer with the color values passed to it. \u201cLevelmap\u201d applies a lookup-table to one or more channels of \u201csource\u201d and puts the results of the lookup into \u201cdestination\u201d. \u201cBlend\u201d blends \u201csource 1\u201d and \u201csource 2\u201d together into \u201cdestination\u201d, based on a ratio passed in. These scripting engine commands, or sometimes referred to as image processing calls or functions, were shipped with many other commands in version 1.0 of FCP from Apple. In addition, there are many scripts written by the users of FCP.\nIn the past, it was assumed that the scripts were written to operate in the RGB color space, 8 bits per color component, and the scripting engine only performed the image processing in the RGB color space."}
-{"text": "There are different solutions for the production of energy from sea waves, for example, solutions based on wave height, wave movement, etc. The present invention utilizes both the wave length and the wave width by the forward wave motion, the backward wave motion and the force of gravity in the most efficient way, something that other proposed solutions do not.\nThe following are examples of the innovation inherent in the present invention as compared with the patents enumerated below, both from the point of view of the operating principle and structure, and from the point of view of energy considerations\u2014allowing greater energy to be attained and produced.    1. (i) U.S. Pat. No. 5,084,630 utilizes only the width of the wave and not the length thereof.            In this patent there is disclosed a system consisting of number of units provided with paddles for each section of wave width. At the bottom of each paddle there are plates. The plates are spaced laterally one from another relative to the direction of wave movement (FIG. 1, the patent top view, shows it). Each paddle activates a pump. In this way, the wave power is not sufficiently utilized because after the wave strikes the paddle and leaves it, it continues to advance and is not longer in contact with the paddle, In FIG. 1, part of the 30 hits the plate 44, leaves the paddle and the wave's part continues to move without utilized its power any more. Other parts of the wave 30 hit the plates 46, 48, 50, 52, 54 and 56, each one wave's part hits only one plate and continue to move without utilizing its power any more. Hence, the wave's power is not fully utilized. The system is disclosed in U.S. Pat. No. 5,084,630 utilizes variations in sea level. For this purpose it is provided with a dedicated subsystem, including a set of separates pumps. Due to this provision the whole system is complicated, as it requires additional control means for coordination of the paddle disposition with the variation in sea level. Furthermore, as the disposition of the paddle is continuously changed and the paddle is not submerged deeply enough to maximally utilize available energy, it is not possible to utilize the energy of those wave's components that move under the wave.        \nThe present invention utilizes:                most efficiently both the length and the width of the wave in the following way:        The lower extremities of the paddles are parallel to the wave width. They are arranged one after other in such a way that from a top view, they are arranged in series. Then, when the wave's part strikes the lower extremity of the first paddle it continues to move, and also strikes the other paddles, until the wave's part is flattened. In this way, the power of the wave is consumed more efficiently so as to get maximum energy from the wave. The energy of the next wave coming after the first wave will be utilized in a similar way.        \n(ii) U.S. Pat. No. 5,084,630 utilizes:                Only the energy in one direction of each paddle. When a paddle, 66+68+70+82+84 as it is shown in FIGS. 3-7, moves pivotal forward, the piston 124+128 as it is shown in FIG. 8, pushes oil within hydraulic cylinder 130 to give the energy to the generator 196, through 188, 146 and accumulartor 140. The oil comes from hydraulic tank 208 through pipe 210. When the paddle returns backward after the wave leaves the paddle, the piston 124+128 returns back and presses some air and/or oil, if desired back to the tank 208 through inlet pipe 212. The accumulator 140 does not get oil and the generator 196 does not get energy. The U.S. Pat. No. 5,084,630 has a one direction piston. The patent uses the paddle energy in one direction, the patent does not utilize the paddle's energy.        \nThe present invention utilizes:                The whole paddle energy. The invention has for each paddle, a bi-directional piston. When a paddle move forward with the wave, the piston inside the hydraulic cylinder, pushes oil to the first accumulator, after the wave leave the paddle, it returns backward and the same piston pushes oil to the second accumulator. The two accumulators send the oil throgh pipes to an hydraulic engine that rotates a generator. The present invention receive energy from both paddle pivotal forward motion and paddle pivotal backward motion.            2. (i) PCT/IL99/00258 or AU PP8932 or GB9800394.0 patent requests (they are the same requests), (hereinafter: \u201cPCT/IL99/00258 patent requests\u201d) show in the application\u2014page 5 paragraphs 4,6, and 7, that the invention has one rod with two pistons with two hydraulic cylinders, both of them send the oil to one accumulator. It is theoretical, in reality the only way that an accumulator can be filled by hydraulic cylinders is, one by one. The hydraulic cylinders above the paddle rod operate one by one and the hydraulic cylinders below the paddle rod operate one by one also, but all the hydraulic cylinders, above and below paddle rod, don't operate one by one. In this way the accumulator can't be filled and work properly. The PCT/IL99/00258 patent requests can be used only to produce energy up to the accumulator. The accumulator should send oil to hydraulic engine that operates the generator. The PCT/IL99/00258 patent requests can't make progress from the accumulator to the generator, because the accumulator does not work. The PCT/IL99/00258 patent requests not can work and produce energy because it uses one accumulator.\nIn the present invention, there are two accumulators (it is shown below in FIG. 4). All pistons oil exit in the first direction, fill the first accumulator with oil one by one, when the paddles move forward, and all pistons oil exit in the second direction, fill the second accumulator with oil one by one when the paddles move backward. This invention can be used to produce electric energy with two accumulators that send the oil to hydraulic engine that rotates the generator.\n(ii) PCT/IL99/00258 patent request claims don't say that in backward paddle unit motion the energy is utilized, i.e. after a wave leave the paddle's plate and the paddle falls pivoting with its gravitation and the backward wave's power, this paddle energy is not utilized. The only energy that is utilized is the energy of a wave that hits the paddle's plates. The gravitation energy of each paddle and the backward wave's energy aren't utilized.\nThe present invention shows in claim 1, that either, in forward paddle unit motion, the paddle's energy is accepted and consumed and, in backward paddle unit motion, the paddle's energy is accepted and consumed.\nThe present invention utilizes:                The whole paddle energy. The invention has a bi-directional piston for each paddle. When a paddle move forward with the wave, the piston inside the hydraulic cylinder, pushes oil to the first accumulator, after the wave leaves the paddle, it returns backward and the same piston pushes oil to the second accumulator. The two accumulators send the oil through pipes to an hydraulic engine that rotates a generator. The present invention receives energy from both motions: paddle pivotal forward motion and paddle pivotal backward motion.        \nThe present invention produces energy in the most efficient way, and the invention's technology is better than PCT/IL99/00258.    3 (i) PCT/IL01/00271 patent request shows in their application different method and different drawings than the present invention, as explained in the PCT/IL01/00271 International Preliminary Examination Report patent request, Item I, in which the FIG. 1-FIG. 6 are the present invention.\nThe following is Item I in the above PCT/IL01/00271 report:                \u201cThe present examination authority considers that the amendments files with the letter dated Jun. 17, 2002 (received on Oct. 21, 2002) introduce subject-matter which extends beyond the content of the application as originally filed, contrary to Article 34(2)(b) PCT. The amendments concerned are the following:        \nWith the above-mentioned letter the application has amended the drawing FIGS. 1 and 2 on drawings page 1/6, and has introduced the drawings FIG. 5 on drawings page 4/6 and drawing FIG. 6 on drawing page 5/6. The amended \u201cconnections\u201d in FIGS. 1 and 2, as well the two additional drawings FIGS. 5 and 6 where not contained in the drawings pages 1/5-5/5 as originally filed.\nIt is not apparent how the amended \u201cconnections\u201d in FIGS. 1 and 2, or the additional FIGS. 5 and 6 are directly and unambiguously derivable from the content of the application as originally filed. Although the respective structural arrangement shown in said FIGS. 1, 2, 5, and 6 may indeed not be excluded from the original subject-matter, it remains that it is not explicitly indicated in the application as originally filed that said respective structural arrangement of FIGS. 5 and 6 necessarily has to be realized in order to perform the claimed invention, or how/why the original \u201cside-to-side connection\u201d of FIGS. 1 and 2 may be superseded by a \u201cbottom-to-bottom connection\u201d.\nHence, the amendments in FIGS. 1 and 2, as well as the additional FIGS. 5 and 6 implicitly introduce subject-matter which extends beyond the content of the application as originally filed, contrary to the requirements of Article 34(2)(b) PCT.\nAs foreseen by rule 70.2(c) PCT, this international preliminary examination report is being established as if the above mentioned amendments and additional had not been made. Hence, it is theoretically assumed for the rest of this report that the drawings figures on files are those shown on drawing page 1/5-5/5 as originally filed. Similarly, the description parts related to the amendments and additions in the drawings are considers not to have made.\nForm PCT/Separate Sheet/409 (Sheet 1) (EPO-April 1997)\u201d\nWhile the connection of the rod 3 is to the side of 14 as in PCT/IL01/00271 patent request and a wave arrives and hits there is an energy loss that comes from the momentum of the 14 side distance multiply by the paddle's weight. When the plant is large each of the paddles have a big weight, which bring to a high momentum with high energy loss.\nFollowing the PCT/IL01/00271 International Preliminary Examination Report, I was forces to ask the PCT patent request for the present invention with the new FIGS. 1 and 2 that I sent to the PCT/IL01/00271. The figures were not published because the PCT/IL01/00271 International Preliminary Examination Report was not published. In addition: please note to the 3rd paragraph above this paragraph, from the words: \u201cAs foreseen by rule 70.2(c) PCT\u201d till the words: \u201cthe drawings are considers not to have made.\u201d\nThis invention bottom-to-bottom connection brings an essential change with high advantage, that the plant energy production is higher than PCT/IL01/00271 patent request.\n(ii) The PCT/IL01/00271 patent request uses, in each one of the paddle unit, as it is written in its claim 1, only one crankshaft mechanism and it can be applied only to very small waves. It can be applied for high waves. When a high wave, (e.g. 5 meters wave height, etc. ) arrives, it hits the paddle plates with high force and one crankshaft mechanism, which is the piston 1 in the present invention drawings, should be very large and the piston size is not realistic piston to be operated.\nThe piston's area is calculated according to the following formula:                The piston area=(The wave's force)/(The pressure in the hydraulic system). Since the piston size is so large and not realistic, the plant must have at least two piston to that each one of the two (for example two) pistons' area is half size of the original one piston and only by splitting the one piston to at least two pistons, the plant can work. The present invention using at least two crankshaft mechanisms, as it is written in the present invention claim 1, which can be applied for high waves that bring high power for high plant electric energy production.        \nEach one of the above two alterations in respect of PCT/IL01/00271 patent request is an essential alteration. According the above 1st alteration the, the present invention plant energy production is much more than in PCT/IL01/00271 patent request, and according to the 2nd alteration, the present invention can work and produce energy with high and low waves, while PCT/IL01/00271 patent request can work only with low waves.\nThe present invention produces energy in the most efficient way, and the invention's technology is better than PCT/IL011/00271 patent request.    4. (i) Patent GB 384603A and Patent JP 57081168 have one crankshaft mechanism for the whole paddles.\nThe energy consideration is:                The power of each wave is known on the basis of the wave's height and width. The wave strikes the first paddle, the second paddle, the third paddle and so forth, until the wave completely flattens, and then the next wave comes along. Just like an accumulator that is depleted, so the system takes all the power minus losses.        \nIn General:\n  Power  =            Energy      Time        =                  The        \u2062                                  \u2062                  plant          '                \u2062        s        \u2062                                  \u2062        power            =                                                                  The                \u2062                                                                  \u2062                total                \u2062                                                                  \u2062                energies                \u2062                                                                  \u2062                of                \u2062                                                                  \u2062                all                \u2062                                                                  \u2062                paddles                \u2062                                                                  \u2062                situated                                                                                        in                \u2062                                                                  \u2062                series                \u2062                                                                  \u2062                between                \u2062                                                                  \u2062                two                \u2062                                                                  \u2062                waves                                                              Time          \u2062                                          \u2062          between          \u2062                                          \u2062          two          \u2062                                          \u2062          waves                    \nThe mathematical principle of calculating the wattage:\nP\u2014The power as calculated according to its height and breadth.\n\u03bc\u2014Paddle efficiency.\nNumber ofWave power before the NthPaddles - NPaddlePaddle Power1PP * \u03bc2P \u2212 P * \u03bc(P \u2212 P * \u03bc) * \u03bc3(P \u2212 P * \u03bc) \u2212 (P \u2212 P * \u03bc) * \u03bc =P(1 \u2212 \u03bc)2 * \u03bcP(1 \u2212 \u03bc)2.........P - The power as calculated according to its height and breadth.\u03bc - Paddle efficiency.\nThus, the series continues, so that the greater the number of paddles, the greater the system power, up to the maximum P.\nThe rod that links the last paddle to the crankshaft mechanism in Patent GB 384603A and Patent JP 57081168 is extremely big and long because of the long distance between two waves.\nIn order to obtain maximum power, it is necessary to provide at least 15 paddles. When the waves are 2 meters to 5 meters height, the distance between two waves is 50-100 meters, so that said rod must also be 50-100 meters!\nBecause such a great rod length is required, the profile of the rod proposed in GB 384603A or in Patent JP 57081168 is extremely thick and large, so that it should not bend under the force of the waves, and quite impractical.\nBecause of the weight of the rod proposed in GB 384603A or in Patent JP 57081168, and the slow wave speed prevailing under said conditions, it is not possible to attain the minimum speed of 500 rpm required to rotate the flywheel of the oil pump (linked to the accumulator that drives an oil engine, which in turn drives a generator and produces the electric energy).\nThus, it should be emphasized, that the arrangements proposed in GB 384603A or in Patent JP 57081168 are completely impractical where waves higher than 2 meters are concerned.\nThus, Patent GB 384603A and Patent JP 57081168 are only applicable when Waves are small and the distance between two waves is short. Where higher waves are concerned, the entire plant proposed by GB 384603 A or by Patent JP 57081168 become dysfunctional and cannot be used.\nThe present invention has a crankshaft mechanism for each paddle unit, without used a rod to link among the paddles. Each paddle unit is operated independently.\nThe present invention is equally practicable for both small and large waves and can be used to obtain high power capacities under both sea conditions, and this is the technological innovation of my patent application over the patents: GB 384603A and JP 57081168.\nThe present invention shows below in claim 1 (the main claim ), that it incorporates the features referring to crankshaft mechanism and a piston and also to providing each of the paddle units with a dedicated crankshaft mechanism to show the difference in patents GB 384603A and JP 57081168 in which the whole plate has one crankshaft mechanism.\n(ii) Patent GB 384603A and Patent JP 57081168 don't utilize the energy of backward paddle unit motion, after a wave leaves the paddle's plate and the paddle falls pivoting with its gravitation and the backward wave's power, this paddle energy is not utilized. The only energy that is utilized is the energy of a wave that hits the paddle's plates. The gravitation energy of each paddle and the backward wave's power, aren't utilized.\nThe present invention utilizes the whole paddle energy.                The invention has a Two-Direction piston for each paddle. When a paddle moves forward with the wave, the piston inside the hydraulic cylinder pushes oil to the first accumulator, after the wave leaves the paddle, it returns backward and the same piston pushes oil to the second accumulator. The two accumulators send the oil through pipes to an hydraulic engine that rotates a generator. The present invention receives energy from both motions: paddle pivotal forward motion and paddle pivotal backward motion.        \nThe present invention shows in claim 1 that either, in forward paddle unit motion, the paddle's energy is accepted and consumed and, in backward paddle unit motion, the paddle's energy is accepted and consumed. This is the second technological innovation of my patent application over patents: GB 384603A and JP 57081168.    5. U.S. Pat. No. 5,244,359 describes a marine energy converter system based on utilization of differences in sea level which it converts into linear movement only of piston 21 which moves within its casing 17. The casing 17 is rigidly connected with pole 3 which is anchored at the bottom of the sea. The above converter can only exploit linear monement, and does not make it possible to exploit the dynamic energy of the sea waves. The above converter is equipped with a floating system piston. The piston is connected to a float 41 and this enables the linear movement of the piston within the piston casing. As a result the piston, within its casing, by virtue of the rising level of water, requires very little force to rise. However, since energy is always a product of force multiplied by distance, the resultant energy is very small as a consequence of the small amount of force exerted.\nThe present invention, however, is based on a different principle, i.e., the exploitation of the dynamic energy of sea waves, which give rise to pivotal movement of the paddles rather than their elevation. The pivotal movement is subsequently converted into linear movement of the piston by means of a crankshaft mechanism 32,40. The piston rises within the casing with relatively great force and the energy produced is thus greater.\nThanks to the structure of the system according to the present invention, it is possible to utilize waves regardless of the height of the wave itself.\nBoth from the point of view of the present invention's operating principle and from the point of view of its particular structure, the U.S. Pat. No. 5,244,359 energy converter is not similar to the present invention.    6. U.S. Pat. No. 4,843,249 describes a hydroelectric system, which utilizes the movement of the waves and converts them into circular movement. Said system includes a turbine 22 whose vanes are moved by the moving wave. The turbine wheel is connected to equipment for the generation of electrical energy.\nThe present invention, by comparison, does not convey circular energy, but rather utilizes the axial movement of the paddles to move the piston in a linear motion.\nAs in the previous case, the above U.S. patent does not void the innovation or the innovation of the present invention, neither in principle nor from the point of view of the structure of the above-described patent.\nIt may also be noted, that the utilization of wave energy in order to create circular energy along, is inferior to the creation of linear motion, on the following counts:\nIn order to achieve better that 90% utilization, for example, a pump activated through circular energy must rotate as speeds of 500-5000 rpm. Below 500 rpm, the efficiency rate is abruptly reduced to about 40%. In order to attain such high speeds, when the wave movement is relatively low, it is necessary to provide additional transmission by means of a few pairs of cog-wheels, thus creating losses in the system. Moreover, if we introduce a number of pairs of cog-wheels, then the total efficiency of the system is reduced, and the energy produced is thus far lower. The cost of the system too, is increased through the provision of may pairs of additional cogwheels.\nHowever, when the piston moves in a linear motion, the immediate efficiency of the system is at least 90% and a linear pump is not limited by any constraints.    7. U.S. Pat. No. 53,311,064 describes a system for generating energy from the movement of sea waves, based on a similar principle to that of the above described hydro-electric system, i.e., the use of a turbine which is rotated by the flow of water. The circular movement is conveyed to the energy generating mechanism by means of a transmission mechanism.\nAll the above given explanations regarding the hydro-electric system according to U.S. Pat. No. 4,843,249, are also valid in this case, both as regards circular motion versus linear motion and the differing structure.    8. Fr Patent 501795 described a system for generating energy from the movement of sea waves.\n(i) The system has one paddle. The low extremity of the paddle's plate is a buoy that moves vertically up and down when the wave reaches the plate. The paddle receives only the height energy from the wave.\nWhen a wave hits a paddle, the paddle goes up in pivotal movement. In this way, each paddle receives energies from: a. the wave height energy which is, potential energy plus b. the wave movement energy, which is kinetic energy. The present invention receives from each paddle much more energy in comparison to the above Fr Patent 501795.\n(ii) The present invention has at least two paddles. The Fr Patent 501795 has one paddle. The advantage of the present invention from this reason is the same as explained in 2(i) above as regards towards U.S. Pat. No. 5,084,630.    9. JP patent 006750 converters can only exploit linear movement, and does not make it possible to exploit the dynamic energy of the sea waves. The hydraulic cylinders 2 and 2a, are pistons connected to a float 1 and this enables the linear movement of the pistons within the piston casing. As a result the pistons within the casing, by virtue of the rising level of water, require very little force to rice. However, since energy is always a product of force multiplied by distance, the resultant energy is very small as a consequence of the small amount of force exerted.\nThe present invention, however, is based on a different principle, i.e., the exploitation of the dynamic energy of sea waves, which give rise to pivotal movement of the paddles rather than their elevation. The pivotal movement is subsequently converted into linear movement of the bi-directional piston by means of a crankshaft mechanism 32, 40. The piston rises within the casing with relatively great force and while the energy produced is a product of force multiplied by distance, the resultant energy is higher as a consequence of the high amount of force exerted. Thus the energy produced according to this system is greater then JP patent 006750.    10. U.S. Pat. No. 4,490,621 described a system for generating energy with two hydraulic motors, one hydraulic motor 121a, for low waves and one hydraulic motor 121b for high waves with very complicated system. When the oil pressure PLa>PLb then first hydraulic motor 121a works and when the oil pressure PLa<PLb, the second hydraulic motor works. The present invention works with all kinds of oil pressure with the same hydraulic motor. If for example, the present invention works with a system of three hydraulic motors, the U.S. Pat. No. 4,490,621 requires 6 hydraulic motors: three for law waves plus three for high waves, which is very expensive and with very complicated system. The present invention is more economical and much simpler than U.S. Pat. No. 4,490,621.    11. U.S. Pat. No. 1,960,622 generate energy does not establish optimum distance of the paddle heights, to make plan maximum energy production, since the rod 20 and the plate 25 move upwards or downwards synchronously, automatically, but not supervised as described in the present invention, with the movement of the sea level at tide time.            The paddles' system according to the present invention stays in the optimum fixed position even when the sea water continues to move upwards and downwards with the tide to create the plant maximum energy production, while in the U.S. Pat. No. 1,960,622, the rod 20 and the plate 25, move upwards or downwards all the time with the sea water at the tide time.        "}
-{"text": "Known porous bodies employed for filter materials or catalytic carriers include those consisting of various materials such as resin, metals or ceramics. Among these, a filter or a catalytic carrier consisting of a ceramics material is generally employed in high temperature or strongly corrosive environments which other materials cannot withstand. A filter or a catalytic carrier consisting of oxide ceramics such as alumina (Al.sub.2 O.sub.3) has already been put into practice.\nAs to a porous body consisting of nonoxide ceramics, on the other hand, only small examples have been put into practice while Japanese Patent Laying-Open No. 63-291882 discloses a silicon nitride based or silicon carbide based porous body prepared by a heat treatment. Further, Japanese Patent Laying-Open No. 1-188479 discloses a method of compacting mixed powder of silicon powder and silicon nitride powder of relatively coarse particles and thereafter nitriding the same thereby preparing a porous body as a solid target.\nAs hereinabove described, it is difficult to use a porous body consisting of resin or a metal in a high temperature or corrosive atmosphere. It is inevitably necessary to employ a porous body made of ceramics for a filter for removing foreign matter from a high-temperature exhaust gas or for a carrier serving as a catalyst for decomposing a harmful matter.\nAs an example of such porous bodies made of ceramics, porous bodies made of alumina have been put into practice. While the porous bodies of alumina are varied in pore size, porosity and bending strength, a porous body having a porosity of 35 to 40% and a mean pore size of 25 to 130 .mu.m has a bending strength of 20 to 35 MPa, whereby the strength of the porous body is insufficient depending on its use.\nIn the silicon nitride based porous body disclosed in the aforementioned Japanese Patent Laying-Open No. 63-291882, the porosity is less than 30% and fluid permeability is insufficient. In general, the strength of ceramics tends to be reduced following an increase in the porosity, and it has been extremely difficult to attain compatibility between the porosity and the strength."}
-{"text": "In recent years, various kinds of electronic devices, such as a tablet terminal, a PDA and a smartphone, have been developed. Most of these electronic devices include touch-screen displays for facilitating input operations by users.\nBy touching a menu or an object, which is displayed on the touch-screen display, by a finger or the like, the user can instruct an electronic device to execute a function which is associated with the menu or object.\nHowever, most of existing electronic devices with touch-screen displays are consumer products which are designed to enhance operability on various media data such as video and music, and are not necessarily suitable for use in a business situation such as a meeting, a business negotiation or product development. Thus, in business situations, paper-based pocket notebooks have still been widely used.\nRecently, a technique for determining a relationship between a tablet terminal and the seating position of a user, based on the direction of writing of handwritten characters, has also been developed.\nIn the meantime, in general, most of tablet terminals have a function of automatically rotating the direction of a screen image in accordance with the direction of the tablet terminal.\nHowever, in a handwriting application which handles a handwritten document, if the direction of the screen image of the application is automatically rotated in accordance with the direction of the tablet terminal, there may be a case in which a feeling of use, such as a feeling of using a real paper-based pocket notebook, cannot be obtained. The reason for this is that if the direction of the screen image is automatically rotated, it would become difficult to handwrite characters, etc. in a free direction on the display screen of the handwriting application.\nOn the other hand, the handwriting application is required to also have a function as a digital tool for handling a handwritten document as digital data."}
-{"text": "1. Field of the Invention\nThe present invention relates to sensors capable of detecting a gas in a fluid and, more particularly, to sensors capable of measuring the presence and the amount of oxygen (O.sub.2).\n2. Description of the Prior Art\nAnalytical methods for the determination of O.sub.2 in a fluid include the amperometric and galvanic methods. These methods are quite rapid, simple in operation, and are specially suited for determining O.sub.2 in either its gaseous state or as it may be found dissolved in a liquid.\nIn the amperometric method, a sensor such as that shown in FIG. 1 is employed. In such a sensor, generally indicated as 10, body 12 is provided having a reservoir 14 therein containing a center post 16 for supporting a gold cathode 18. A membrane holder 20, cap 22, and O-ring 23 are provided to hold a polymeric membrane material 24, such as polytetrafluoroethylene of a thickness of from 0.00025 inches to 0.002 inches, stretched over the gold cathode thereby completely enclosing reservoir 14. Reservoir 14 is then filled with an electrolyte 26, typically a 5% KCl solution, either buffered or unbuffered. When the sensor is to be used in an application where temperatures in excess of 100.degree. F. are generally encountered, the electrolyte is generally mixed with typical antifreeze compounds such as ethylene glycol. Additionally, an anode 28, typically silver, is disposed within reservoir 14 in contact with electrolyte 26. A potential of about 750 millivolts is applied between the anode 28 and the cathode 18 by means 29 connected thereto. In stretching the membrane 24 over cathode 18, a very minimal amount of electrolyte 26 is contained between gold cathode 18 and membrane 24. As sample fluid is brought in contact with membrane 24, O.sub.2 diffuses through membrane 24 to contact the gold cathode 18 in the presence of the electrolyte 26. A current flow results which is linear with the partial pressure of O.sub.2 being sampled. Thus, this current can be measured and correlated to the amount of O.sub.2 in the sample.\nIn the galvanic method, a sensor such as that shown in FIG. 2 is employed, wherein corresponding numbers indicate corresponding parts. The sensor employed in the galvanic method typically has a silver cathode 18', a lead anode 28' and a readout device 29', e.g., an ammeter.\nWhile these methods are satisfactory in many applications, they suffer drastically in atmospheres of high carbon dioxide (CO.sub.2) such as encountered in monitoring automobile exhaust or stack gas emissions. Such sensors as that of FIG. 1 and FIG. 2, as described above, become less sensitive to O.sub.2 upon even brief exposure to high concentrations of CO.sub.2 and may take several hours to recover so as to indicate the proper value of O.sub.2. The response of such sensors is highly dependent on the pH of the electrolyte at the interface between the gas diffusing through the membrane and the cathode. At low, or acid, pH levels the response is low. At high, or base, pH levels the response is higher. Typical responses by a prior art sensor are shown in FIG. 3 and FIG. 4. In FIG. 3, the sensor was first exposed to ambient air containing approximately 21% O.sub.2. It was then exposed to pure nitrogen (N.sub.2). FIG. 3 shows the response which was both expected and achieved. Upon exposure to N.sub.2, the response dropped to the zero line. Upon exposure to air, the response climbed to the level indicating approximately 21% O.sub.2. This cycle was repeatable without problem. Referring now to FIG. 4, the expected and actual response of a prior art sensor is shown when the sensor was exposed to ambient air and then exposed to a mixture of 15% CO.sub.2 plus 3% O.sub.2 and the balance nitrogen. When exposed to the mixture, the expected response is for the output to drop to the 3% level, being an indication of the 3% O.sub.2 content of the mixture. Upon exposure to ambient air, it is expected that the output will climb to the 21% O.sub.2 level of the ambient air. The actual response, however, was not as anticipated. When the sensor was exposed to the mixture, the response fell to the expected 3% O.sub.2 level. When the sensor was subsequently exposed once again to the ambient air sample, the output overshot the 21% level, then reversed and undershot the 21% level, and then slowly approached the 21% level asymptotically. It was found that the recovery period required for the output to attain the actual 21% level varied depending both on the duration of exposure to CO.sub.2 and the amount of CO.sub.2 in the sample. For example, when using such a prior art sensor on an automobile exhaust, an exposure for a period one minute to the exhaust gases resulted in a recovery period on the order of two to three hours before an accurate ambient response could be attained.\nThis phenomenon is a result of the small volume of electrolyte trapped adjacent to the gold cathode by the membrane. This is typically on the order of 1 microliter. As previously mentioned, the response of such a cell is dependent on the pH of the electrolyte. When CO.sub.2 is introduced, carbonic acid is formed which, when mixed with such an extremely small volume of basic electrolyte, results in a change of the pH of the electrolyte adjacent the gold cathode. In the typical electrolyte having a pH of approximately 13.5, the introduction of carbonic acid having a pH in the order of 4.5 results in a change of pH of the electrolyte trapped adjacent the gold cathode to a level of approximately 9. The response of the electrode will be correspondingly reduced until such time as the pH can attain its normal value by diffusion of normal electrolyte into the space between the membrane. The initial overshoot observed is, presumably, caused by the sudden change in pH and the unsettling of the electrolyte in the cathode area.\nAttempts at improving the performance of O.sub.2 sensors are not new in the art. It is well known that changes in pH of the electrolyte adjacent the cathode will change the response of the electrode. On the other hand, it is known that the thickness of the electrolyte in this same area affects the sensitivity of the electrode to oxygen. Thus, ideally, the spacing between the membrane and the cathode is kept minimal while means are provided for allowing the free movement of the electrolyte through the space. Thus, in the prior art, it has been suggested to roughen the surface of the gold cathode, provide channels therein for the movement of electrolyte, and depose porous materials between the membrane and the cathode to provide channels for the movement of the electrolyte.\nSuch prior art suggestions have resulted in O.sub.2 sensors of marginal sensitivity and poor response times for certain applications. In particular, automobile exhaust analysis and flue gas analysis provide environments imposing restrictions beyond the capabilities of prior art O.sub.2 sensors employing such techniques. In the field of automobile exhaust gas analysis, the ability to cycle and recover at rapid rates is imperative in \"assembly line\" type testing environments.\nAccordingly, it would be highly desirable to have an electrochemical O.sub.2 sensor capable of yielding accurate and reproducible date in applications wherein high CO.sub.2 levels are encountered at both ambient, or low, and elevated temperatures. Such a sensor would be highly useful for monitoring the O.sub.2 content of automobile exhaust and stack gas emissions."}
-{"text": "Display devices are used for presentation of video to a viewer. The video to be presented may be transmitted to the display device from a source unit in the form of a video signal. A display device receiving a video signal will render video output corresponding to data conveyed by the video signal. However, if the display device improperly interprets the video signal, which could occur, e.g., if the video signal includes excessive noise, the corresponding video output may be disrupted or compromised."}
-{"text": "1. Field of the Invention\nThe present invention relates to a matrix type display apparatus such as a liquid crystal display (hereinafter referred to as an LCD), and particularly to a technique for the enlarged display of a screen.\n2. Description of the Related Art\nFIG. 2 is a structural diagram of a conventional LCD.\nThis LCD displays image signals supplied from a personal computer or the like (referred to below as a PC) on a color screen having a size of 1024 pixels horizontally by 768 pixels vertically. This LCD comprises a control circuit 10, a display signal circuit 20, a scan signal circuit 30, and a liquid crystal panel 40.\nRed, green, and blue color image signals formed from R, G, and B signals, clock signals CLK showing timings of a sampling of these R, G, and B signals, horizontal synchronizing signals HSYN, and vertical synchronizing signals VSYN are fed to the input side of the control circuit 10 from an unillustrated PC. On the basis of these signals, the control circuit 10 then generates start signals EI, clock signals CK, R, G, B data, and strobe signals STB and outputs these to the display signal circuit 20. The control circuit 10 also has the function of generating start signals ST and clock signals CP and outputting these to the scan signal circuit 30.\nThe display signal circuit 20 is provided with a 1024 stage shift register 21 which corresponds to the 1024 pixels in the horizontal direction, data latches 22 and 23, and a display drive section 24. The shift register 21 sequentially shifts and holds the start signals EI based on the clock signals CK and the held contents of each stage are fed to each stage of the data latch 22 as latch signals L1, L2, etc up to L1024. The data latch 22 holds the R, G, B data fed from the control circuit 10 based on the latch signals L1 to L1024.\nThe data latch 23 stores the R, G, B data of the 1024 horizontal pixels held by the data latch 22 at the strobe signal STB timing. The display drive section 24 generates display voltages S1, S2 etc up to S3072 which correspond to R, G, B data of the 1024 pixels stored in the data latch 23 and outputs these display voltages.\nThe scan signal circuit 30 is provided with a 768 stage shift register 31 and scan drive section 32 for sequentially scanning and displaying the 768 pixels in the vertical direction in horizontal line units. The shift register 31 holds the start signals ST in the initial stage register as scan signals and then sequentially shifts these scan signals in accordance with the timing of the clock signals CP. The contents at each stage of the shift register 31 are fed to the scan drive section 32 and are output as scan voltages G1, G2, etc up to G768.\nThe liquid crystal panel 40 is provided with 768 X electrodes X1, X2, etc up to X768 arranged at equal spacing in a line direction, and 1024 groups (note that each group is made up of 3 electrodes corresponding to R, G, and B) of Y electrodes Y1, Y2, etc up to Y3072 arranged at equal spacing in a column direction. Color pixel display is performed in line units in accordance with display voltages S1 to S3072 applied respectively to the Y electrodes Y1 to Y3072 at locations where an X electrode Xj to which a scan voltage Gj (wherein j is a number from 1 to 768) is applied intersects with a Y electrode Y1 to Y3072.\nIn this type of LCD, when R, G, B signals, clock signals CLK, horizontal synchronizing signals HSYN, and vertical synchronizing signals VSYN are supplied from a PC, horizontal cycle start signals EI and, thereafter, in synchronization with the clock signals CLK, R, G, B data is sequentially transmitted in single pixel units from the control circuit 10 to the display signal circuit 20. In the shift register 21, latch signals L1 to L1024 are generated in accordance with the clock signals CK. R, G, B data is further sequentially held in the data latch 22 in synchronization with the latch signals L1 to L1024.\nWhen the R, G, B data of the 1024 pixels of a single line is held in the data latch 22, a strobe signal STB is output from the control circuit 10. Accordingly, the R, G, B data of the single line in the data latch 22 is also stored in the data latch 23. In the display drive section 24, display voltages S1 to S3072 are generated based on the R, G, B data of the 1024 pixels stored in the data latch 23 and are output.\nMoreover, start signals ST for each vertical cycle and clock signals CP for shifting the start signals ST and generating scan signals are output from the control circuit 10 to the scan signal circuit 30. When the start signals ST are supplied, the output signal of the initial stage of the shift register 31 of the scan signal circuit 30 is set to a level xe2x80x9cHxe2x80x9d for designating display and output signals of subsequent stages are set to a level xe2x80x9cLxe2x80x9d for designating non-display. Output signals of each stage of the shift register 31 are then sequentially shifted one stage at a time towards the rear in synchronization with the clock signals CP which correspond to the horizontal cycle. The output signals of each stage of the shift register 31 are fed to the scan drive section 32 and scan voltages G1 to G768 are generated for each line in accordance with the display/non-display and output.\nThe display voltages S1 to S3072 which correspond to the R, G, B data of a single line output from the display drive section 24 are fed to the Y electrodes Y1 to Y3072 of the liquid crystal panel 40. Moreover, the scan voltages G1 to G768 output from the scan drive section 32 are fed to the X electrodes X1 to X768 of the liquid crystal panel 40. The timing of the strobe signals STB fed from the control circuit 10 to the display signal circuit 20 is substantially identical to the timing of the clock signals CP fed to the scan signal circuit 30. Therefore, a single line which corresponds to the X electrode Xj driven by the scan voltage Gj output from the scan drive section 32 is displayed through the display voltages S1 to S3072 based on the R, G, B data for the single line stored in the data latch 23. In addition, all other X electrodes other than the X electrode Xj are placed in a non-display state.\nA screen is displayed by the X electrodes X1 to X768 being sequentially driven from top to bottom by the scan voltages G1 to G768 sequentially output from the scan drive section 32.\nHowever, this type of conventional LCD has the following problems.\nFor example, when an image signal supplied from a PC is for displaying on a 640 horizontal pixel by 480 vertical pixel screen, no pixel data exists for displaying at the right hand side and bottom of a liquid crystal panel 40 having 1024 horizontal pixels by 768 vertical pixels. Therefore, the problem has existed that the displayed screen is small and the display position is offset to the top left.\nIn order to solve this problem, attempts have been made to provide a processing device and frame memory corresponding to the resolution of the liquid crystal panel 40 in the LCD. Image signals fed from the PC are interpolated by this processing device and converted into image data for a 1024 horizontal pixel by 768 vertical pixel screen and displayed.\nHowever, in this type of method, because a large amount of frame memory and a high speed processing device are necessary, the problem arises that costs are increased and the amount of power needed to run the processing device at high speed also increases.\nThe present invention solves the above problems in the conventional technology and provides a matrix type display apparatus such as an LCD which has a simple circuit structure and which is capable of enlarging the display of a screen.\nIn order to solve the above problems, the first aspect of the present invention is a matrix type display apparatus comprising: display means having M number of X electrodes (wherein M is a number greater than one) arranged in parallel, N number of Y electrodes (wherein N is a number greater than one) arranged so as to intersect the X electrodes, and display elements provided at each point of intersection of the X electrodes and Y electrodes, the display means performing a matrix type display at the points of intersection of the N number of Y electrodes and the X electrodes which are driven by scan voltage in accordance with a display voltage applied to each Y electrode; scan drive means for generating the scan voltage in order to sequentially drive the X electrodes; and display drive means for holding image data which corresponds to the X electrodes driven by the scan drive means, and for generating the display voltage based on the held image data and driving each of the Y electrodes, wherein the scan drive means and the display drive means have the following structure.\nNamely, the scan drive means is provided with M number of flip-flop circuits (referred to below as xe2x80x9cFF circuitsxe2x80x9d) and switches for altering the connections of the M number of FF circuits and, in normal display mode, the scan drive means forms shift registers having M number of stages by connecting the M number of FF circuits in series and also generates the scan voltages based on each output signal of the M number of FF circuits and, in enlarged display mode, the scan drive means forms shift registers having m number of stages (wherein m is a number less than M) by connecting a portion of the M number of FF circuits in parallel at a constant ratio and also generates the scan voltages based on each output signal of the M number of FF circuits.\nMoreover, the display drive means is provided with N number of FF circuits and switches for altering the connections of the N number of FF circuits and, in the normal display mode, the display drive means forms shift registers having N number of stages by connecting the N number of FF circuits in series and also holds the image data based on each output signal of the N number of FF circuits and generates the display voltages in accordance with the image data and, in the enlarged display mode, the display drive means forms shift registers having n number of stages (wherein n is a number less than N) by connecting a portion of the N number of FF circuits in parallel at the constant ratio and also holds the image data based on each output signal of the N number of FF circuits and generates the display voltages in accordance with the image data.\nThe second aspect of the present invention is the matrix type display apparatus according to the first aspect, wherein switches in the scan drive means and the display drive means are formed from analog switches (referred to below as xe2x80x9cSWxe2x80x9d) controlled by mode signals which designate the normal display mode or the enlarged display mode.\nAccording to the first and second aspects of the present invention, because the above described structure is employed, the following operation can be performed.\nIn normal display mode, the M number of FF circuits in the scan drive means are connected in series by switches such as SW or the like. Scan voltages are sequentially generated from each FF circuits and fed to the X electrodes of the display means. The N number of FF circuits in the display drive means are also connected in series by switches. Image data based on output signals sequentially output from each FF circuit is held, and display voltages are generated in accordance with the image data and fed to the Y electrodes of the display means.\nIn contrast, in enlarged display mode, the M number of FF circuits in the scan drive means are connected in parallel and in series at a constant ratio by switches. Scan voltages are generated based on output signals from each FF circuit and fed to the X electrodes of the display means. The N number of FF circuits in the display drive means are also connected in parallel and in series at a constant ratio by switches. Image data based on output signals from each FF circuit is held, and display voltages are generated in accordance with the image data and fed to the Y electrodes of the display means.\nThe third aspect of the present invention is the matrix type display apparatus according to the first and second aspects, wherein the scan drive means is provided with a plurality of integrated circuits used for shifting which are connected in a cascade connection and, in the enlarged display mode, the scan drive means simultaneously drives outputs of scan voltages for front or rear non-display regions based on selection signals which correspond to the integrated circuit connection position.\nThe fourth aspect of the present invention is the matrix type display apparatus according to the third aspect, wherein selection signals for a plurality of integrated circuits in the scan drive means are directly fed from ground voltage or power supply voltage supplying the integrated circuits.\nAccording to the third and fourth aspects of the present invention, the following operation is performed in the scan drive means in the matrix type display apparatus according to the first and second aspects of the present invention.\nSelection signals corresponding to the connection positions at which a plurality of integrated circuits used for shifting are connected in a cascade connection in the scan drive means are directly fed, for example, from a power supply and ground voltage identical to the integrated circuits. In enlarged display mode, the output of scan voltages for front and rear non-display regions is stopped based on the selection signals."}
-{"text": "1. Field of the invention\nThis invention relates in general to wellhead equipment for oil and gas wells, and in particular to an emergency casing hanger seal.\n2. Description of the Prior Art\nIn a well of the type concerned herein, a wellhead housing has a landing profile or shoulder within its bore. When running casing conventionally, a casing hanger is installed on the upper end of the string of casing. The casing hanger lands on the landing shoulder in the bore of the wellhead housing.\nAfter cementing, a seal is positioned between the casing and the wellhead housing. The seal locates between machined surfaces on the wellhead housing and the casing hanger. A tubing hanger may be installed over the uppermost casing hanger. The tubing hanger is normally held by lockdown screws if the wellhead housing is located at the surface. The tubing hanger secures to the tubing which extends into the well. A tubing hanger seal locates between the tubing hanger and the wellhead housing.\nOccasionally, the casing will not proceed smoothly to the bottom of the well. When this occurs, the casing hanger will not be properly positioned to land in the wellhead housing. Generally, when this happens, the casing cannot be retrieved to the surface and becomes stuck. In that case, the casing must be cut above the landing shoulder after cementing. In the prior art, the casing is supported by slips in the wellhead housing.\nA problem exists in sealing against the casing stub, because the casing does not have a smooth machined surface for receiving a seal. The casing outer diameter has a high dimensional variation. The outer diameter may be slightly oval shaped. The surface of the casing may have many defects, such as rust, pock marks and tong marks. Various seals have proposed for sealing against the casing stub. However, improvements in locking the seal in the wellhead housing are desired.\nAnother desirable feature would be a means that would prevent the casing from moving upward due to temperature increase in the well during production. Movement of the casing could have damaging effects on the seal. Also, it would be desired to have a tubing hanger lockdown that did not employ lockdown screws, and would allow the seal to be removed without removing the lockdown."}
-{"text": "1. Field of the Invention\nAbsorption refrigeration systems having means for automatically controlling the energizing medium.\n2. Description of the Prior Art\nIn U.S. Pat. No. 3,895,499 issued to N. E. Hopkins on July 22, 1975, there is described an absorption refrigeration system in which the cooling tower water is essentially uncontrolled and allowed to vary as influenced by the temperature of the ambient air flowing through the cooling tower in contact with the water spray devices. In the Hopkins system, the heat input to the generator is controlled in response to the entering or leaving chilled water temperature as is conventional in absorption machines. The basic concept embodied in the Hopkins system is a means for transferring solution from one point to another to prevent a dangerous operating condition.\nIn U.S. Pat. No. 4,164,128 issued to A. B. Newton on Aug. 14, 1979, there is described an absorption refrigeration system which is adapted to be powered by solar energy or a source of relatively low temperature waste heat. A control system automatically adjusts the maximum energizing temperature as the temperature of the cooling water and the temperature of chilled water supplied to the load vary from time to time.\nAn article by A. B. Newton entitled \"Optimizing Solar Cooling Systems\", ASHRAE Journal, November 1976, pp. 26-31, suggests that plural compartments for the storage of heated medium at different temperatures may be used in a solar powered absorption system. No specific means are described for controlling the selection of the storage compartments."}
-{"text": "Providing defense mechanisms against computer hackers and attacks on computing infrastructure has been an increasingly significant concern with the growing importance of network-based computing. With the growth of network-based computing over the Internet and the growing reliance on cloud computing services, protecting computers and servers connected to a network from malicious network-based attacks has become a significant issue for businesses and organizations. For example, hackers may attack a computing service over a network using denial-of-service attacks. SYN flooding is a particularly popular denial-of-service tactic. In a SYN flood attack, an attacker sends a succession of SYN requests to a target system in an attempt to overwhelm the target system's resources so that it becomes unresponsive to legitimate requests. Computing service providers must account for and protect against such attacks while maintaining an effective level of service to their customers."}
-{"text": "1. Field of the Invention\nThe present invention relates to a data recording disk, a disk drive device and a method for writing servo patterns on a data recording surface of a rotating data recording disk, and more particularly to a method, data recording disk and disk drive having servo patterns that improve efficiency of head positioning and data storage.\n2. Description of Related Art\nA giant magneto resistive (GMR) head has been developed as a next generation read/write head of a hard disk drive device. The GMR head can generate an output signal which is larger than an output signal generated by a prior magneto resistive (MR) head. The capability of generating the large output signal causes a reduction of a head width of the GMR head.\nThe GMR head is fabricated by the fabrication process of the MR head, and hence a head width of the GMR head is determined by a tolerance of the fabrication process of the MR head. It is assumed that the head width of the prior MR head is 2 xcexcm, (1 xcexcm=10xe2x88x926 m), the head width of the GMR head is 1 xcexcm, and the tolerance of the fabrication technology of the MR head is 0.4 xcexcm. Since the tolerance of the 0.4 xcexcm is applied to both the MR head and GMR head, the head width of the MR head is 2 xcexcmxc2x10.4 xcexcm, that is, the track width of the MR head varies from 1.6 xcexcm to 2.4 xcexcm, and a variance of the head width is 20%. The head width of the GMR head is 1 xcexcmxc2x10.4 xcexcm, that is, the head width of the GMR head varies from 0.6 xcexcm to 1.4 xcexcm, and a variance of the head width is 40%. In this manner, the variance of the head width of the MR head is 20%, while the variance of the head width of the GMR head is 40% which is larger than that of the MR head.\nAn uniform track pitch is used in the track arrangement of the prior MR technology. A track pitch is a distance between a track center of one track and a track center of the next track.\nIf the uniform track pitch of the MR head is used for the GMR head, the following first problem, as shown in FIG. 1 occurs. As described above, the head width of the GMR head varies from 0.6 xcexcm to 1.4 xcexcm. For simplifying the description, three kinds of the GMR heads 1, 2 and 3 are shown in FIG. 1, i.e., the GMR head 1 has the 1.4 xcexcm head width, the GMR head 2 has the 1.0 xcexcm head width, and the GMR head 3 has the 0.6 xcexcm head width. The track width of all the tracks 1, 2 and 3 are the same width, i.e., 1.4 xcexcm, which is the largest width of the GMR head, and the uniform track pitch is used for all the tracks. A gap is provided for separating the adjacent two tracks. And, the center of the GMR heads 1, 2 and 3 are aligned to the track center of each track, respectively. It is apparent that useless or wasted space exists between the GMR heads 1 and 2, and between the GMR heads 2 and 3.\nThe following second problem occurs with respect to the servo patterns provided in the track. The servo patterns A and B are written to define the tracks, and read by the head to generate a feedback signal which is supplied to a hard disk control device. The detail of the servo patterns A and B are shown in FIG. 4(B) and (C). The servo pattern A has alternately arranged opposite magnetization direction, and the servo pattern B has an uniform magnetization direction. The hard disk control device responds to the feedback signal to position the center of the head on the track center. It is assumed that the GMR heads 1, 2 and 3 are shifted from the track center by a distance of 10% of the track pitch, as shown by the dashed line in FIG. 1. For example, it is assumed that the track pitch=track widthxc3x971.1=1.4 xcexcm=1.1xc3x971.6 xcexcm, and the shift distance=track pitchxc3x970.1=0.16 xcexcm. In this case, the GMR heads 1, 2 and 3 are shifted with respect to the servo pattern B, and the variation of the levels of the feedback signals sensed by the GMR heads 1, 2 and 3 are as follows.\nIn the case of the GMR head 1 of the 1.4 xcexcm head width: 0.16 xcexcm/1.41 xcexcm=0.114. That is, the level of feedback signal is reduced by 11.4%. In the case of the GMR head 2 of the 1.01 xcexcm head width: 0.16 xcexcm/1.0 xcexcm=0.16. That is, the level of feedback signal is reduced by 16%. In the case of the GMR head 3 of the 0.6 xcexcm head width: 0.16 xcexcm/0.6 xcexcm=0.267. That is, the level of feedback signal is reduced by 26.7%.\nThe variation of the levels of the feedback signal sensed by the prior MR head of 2.4 xcexcm head width, the prior MR head of 2.0 xcexcm head width, and the prior MR head of 1.6 xcexcm head width are, as follows. It is assumed that the track pitch=track widthxc3x971.1=2.4 xcexcmxc3x971.1=2.6 xcexcm, and the shift distance=track pitchxc3x970.1=0.26 xcexcm.\nIn the case of the MR head of the 2.4 xcexcm head width: 0.26 xcexc m/2.4 xcexc m=0.108. That is, the level of the feedback signal is reduced by 10.8%. In the case of the MR head of the 2.0 xcexcm head width: 0.26 xcexcm /2.0 xcexcm=0.13. That is, the level of the feedback signal is reduced by 13%. In the case of the MR head of the 1.6 xcexcm head width: 0.26 xcexcm/1.6 xcexcm=0.163. That is, the level of the feedback signal is reduced by 16.3%.\nIn this manner, when the head is shifted by the distance of track pitch multiplied by A, the level of the feedback signal of the MR head to the hard disk control device varies from 10.8% to 16.3%, while the level of the feedback signal of the GMR head to the hard disk control device varies from 11.4% to 26.7%.\nThe range of the variation of the feedback signals to the hard disk control device should be small, since if the range of the variation becomes large, undesirable phenomenon, such as an oscillation of the entire servo system may occur.\nIt can be seen then that there is a need for a method, data recording disk and disk drive that eliminates useless or wasted space on a data recording disk.\nIt can also be seen that there is a need for a method, data recording disk and disk drive that provides a track pitch that varies depending o the track width.\nIt can also be seen that there is a need for a method, data recording disk and disk drive that provides a uniform level of variation I the feedback signals.\nTo overcome the limitations in the prior art described above, and to overcome other limitations that will become apparent upon reading and understanding the present specification, the present invention discloses a data recording disk that includes a first data recording surface and a second data recording surface, wherein a track pitch of adjacent data recording tracks of the first data recording surface differs from a track pitch of adjacent data recording tracks of the second data recording surface.\nA data recording disk according to the present invention includes a first data recording surface and a second data recording surface, wherein the number of data recording tracks of the first data recording surface differs from the number of data recording tracks of the second data recording surface, and the a gap width between the data recording tracks of the first data recording surface differs from a gap width between the data recording tracks of the second data recording surface.\nA data recording disk according to the present invention includes a first data recording surface and a second data recording surface, wherein the number of data recording tracks of the first data recording surface differs from the number of data recording tracks of the second data recording surface.\nA disk drive device according to the present invention includes data recording surfaces on a rotating data recording disk, each of the data recording surfaces is provided with a read/write head, wherein the number of data recording tracks of a first data recording surface differs from the number of data recording tracks of a second data recording surface, and a gap width between the data recording tracks of the first data recording surface differs from a gap width between the data recording tracks of the second data recording surface.\nAll data recording tracks of the first data recording surface have substantially the same track width, and all gaps of the first data recording surface have substantially the same gap width.\nAll data recording tracks of the second data recording surface have substantially the same track width, and all gaps of the second data recording surface have substantially the same gap width.\nData is written in a direction from an outer most data recording track to an inner most data recording track of the first data recording surface, then the data is written in a direction from an outer most data recording track to an inner most data recording track of the second data recording surface.\nAddresses of data recording tracks are successively assigned from an outer most data recording track to an inner most data recording track of the first data recording surface, and the address of the inner most data recording track of the first data recording surface is followed by an address of an outer most data recording track of the second data recording surface, and the addresses on the second data recording surface are successively assigned from the outer most data recording track to an inner most data recording track of the second data recording surface.\nA method according to the present invention for writing servo patterns for defining a plurality of data recording tracks on a rotating data recording disk by a read/write head contained in a disk drive device, includes positioning the read/write head at one position along a radial direction of the rotating data recording disk, writing servo patterns on the rotating data recording disk by the read/write head positioned at the one position, reading the servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50xe2x88x92Gw/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Gw is a predetermined gap width between the data recording tracks, stopping the read/write head when the detection is made, and overlappingly writing the servo patterns on the rotating data recording disk by the stopped read/write head.\nA method according to the present invention for writing servo patterns for defining a plurality of data recording tracks on a rotating data recording disk by a read/write head contained in a disk drive device, includes specifying a gap width between adjacent data recording tracks, writing servo patterns for defining a plurality of data recording tracks separated by the specified gap width on the rotating data recording disk by the read/write head, detecting that the total number of data recording tracks is smaller than a predetermined number, specifying an overlap width of adjacent data recording tracks, and rewriting servo patterns for defining a plurality of data recording tracks overlapped each other by the overlap width on the rotating data recording disk by the read/write head, the rewriting of servo patterns including positioning the read/write head at one position along a radial direction of the rotating data recording disk, rewriting the servo patterns on the rotating data recording disk by the read/write head positioned at the one position, reading the rewritten servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the rewritten servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50+Y/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Y is the overlap width, stopping the movement of the read/write head when the detection is made, and rewriting the servo patterns on the rotating data recording disk by the stopped read/write head stopped.\nA method according to the present invention for writing servo patterns for defining a plurality of data recording tracks on each of a plurality of data recording surfaces on rotating data recording disks by a read/write head provided on each of the plurality of data recording surfaces contained is a disk drive device, includes specifying a gap width between adjacent data recording tracks, writing servo patterns on a plurality of servo tracks for defining the plurality of data recording tracks separated by the specified gap width on one data recording surface by its read/write head, storing a count value indicating the number of servo tracks written on the one data recording surface, repeating the writing of servo patterns and storing the count value until the servo patterns are written on all data recording surfaces, determining whether the total number of data recording tracks written on the data recording surfaces exceeds a predetermined number, or not, specifying an overlap width of adjacent data recording tracks, if the total number of data recording tracks written on data surfaces does not exceed the over the top width, finding the data recording surface containing the least number of data recording tracks by determining the count values of the all data recording surfaces, and rewriting the servo patterns for defining the plurality of data recording tracks overlapped each other by the overlap width on the data recording surface found to contain the least number of data recording tracks to its read/write head, the rewriting including positioning the read/write head at one position along a radial direction of the data recording surface, rewriting servo patterns on the data recording surface by the read/write head positioned at the one position, reading the rewritten servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the rewritten servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50+Y/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Y is the overlap width, stopping the movement of the read/write head when the detection is made, and rewriting the servo patterns on the data recording surface by the stopped read/write head stopped.\nA method according to the present invention for writing servo patterns for defining a plurality of data recording tracks on each of a plurality of data recording surfaces on rotating data recording disks by a read/write head provided on each of the plurality of data recording surfaces contained is a disk drive device, includes specifying a gap width between adjacent data recording tracks, positioning the read/write head at one position along a radial direction of one data recording surface, setting a count value of a counter to an initial value, writing servo patterns on the data recording surface by the read/write head positioned at the one position, reading the servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50xe2x88x92Gw/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Gw is the specified gap width between the data recording tracks, stopping the read/write head when the detection is made, incrementing the count value, overlappingly writing the servo patterns on the data recording surface by the stopped read/write head, repeating the reading, moving, reading and stopping until the servo patterns are entirely written on the one data recording surface, electing next data recording surface and repeating the positioning, setting, writing, reading, moving, reading, stopping, incrementing, overlapping and repeating, determining whether the total number of data recording tracks written on all data recording surfaces exceeds a predetermined number, or not, specifying an overlap width of adjacent data recording tracks, if the total number of data recording tracks written on all data recording surfaces does not exceed the predetermined number, finding the data recording surface containing the least number of data recording tracks by determining the count values of the all data recording surfaces, and rewriting the servo patterns for defining the plurality of data recording tracks overlapped each other by the overlap width on the data recording surface found containing the least number of data recording tracks by its read/write head, the rewriting including positioning the read/write head at one position along a radial direction of the data recording surface, rewriting servo patterns on the data recording surface by the read/write head positioned at the one position, reading the rewritten servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the rewritten servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50+Y/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Y is the overlap width, stopping the movement of the read/write head when the detection is made, and rewriting the servo patterns on the data recording surface by the stopped read/write head stopped.\nThe total number of data recording tracks written on the all data recording surfaces is determined by referring to the count values of the all data recording surfaces.\nThese and various other advantages and features of novelty which characterize the invention are pointed out with particularity in the claims annexed hereto and form a part hereof. However, for a better understanding of the invention, its advantages, and the objects obtained by its use, reference should be made to the drawings which form a further part hereof, and to accompanying descriptive matter, in which there are illustrated and described specific examples of an apparatus in accordance with the invention."}
-{"text": "1. Field of the Invention\nThe present invention relates to an optical recording medium, in particular to a write once optical recording medium capable of recording information therein and reproducing the recorded information therefrom by causing the optical characteristics, such as transmittance and reflectance, of a recording material for use in the optical recording medium to change by the application of an optical beam thereto. In addition, the present invention also relates to a method of recording information in the above-mentioned optical recording medium.\n2. Discussion of Background\nThe wavelength of a laser beam used for currently employed write once optical disk systems of a write once read many type (WORM), and a compact disc-recordable (CD-R) type is in the range of 770 nm to 790 nm, and a recording medium for the write once optical disk system is designed in such a fashion that recording and reproduction can be performed in the above-mentioned wavelength range.\nSooner or later, it will be inevitable that the capacity of such a recording medium must be significantly increased in accordance with the amount of information to be handled. It is also inevitable that the wavelength of the laser beam used for the recording and reproduction will be shortened.\nIn recent years, there have been proposed many types of write once optical disks for data recording and reproduction, using a cyanine dye or a phthalocyanine dye as a recording material.\nFor example, the write once optical disks using as the recording material a cyanine dye are disclosed in Japanese Laid-Open Patent Applications 57-62093, 58-56892, 58-112790, 58-114989, 59-85791, 60-83236, 60-89842, and 61-25886. Furthermore, the write once optical disks using as the recording material a phthalocyanine dye are disclosed in Japanese Laid-Open Patent Applications 61-150243, 61-177287, 61-154888, 61-246091, 62-39286, 63-37791, and 63-39888.\nHowever, there has not yet been developed any recording material that has high light resistance and excellent preservation stability, and is capable of performing the recording and reproducing operation using an optical pick-up of a laser beam with a wavelength of 700 nm or less.\nThe currently employed CD-R is also designed in such a fashion that recording and reproduction can be performed by a laser beam with a wavelength in the range of 770 nm to 790 nm.\nIn the case of the above-mentioned CD-R system, the increase in capacity of the recording medium also becomes inevitable, and there is the trend toward shortening of the wavelength of the laser beam to be used for data recording and reproduction.\nIn CDs and CD-ROMs which are currently employed, aluminum (Al) is coated on the concave and convex portions of the substrate therefor and the wavelength-dependence of the reflectance of Al is so small that even if the wavelength of the laser beam used therefor is shortened in the future, reproduction of information recorded in such CDs and CD-ROMs is possible.\nIn contrast to this, the recording layer of the currently employed CD-R, which contains therein a dye with a maximum absorption wavelength in the range of 680 nm to 750 nm, is designed so as to exhibit a maximum reflectance when a laser beam with a wavelength of 770 nm to 790 nm is applied thereto, with the optical characteristics and constants and the thickness and structure of the recording layer being taken into consideration. Therefore, when a laser beam with a wavelength of 700 nm or less is applied to the currently employed CD-R, the reflectance exhibited by the recording layer is extremely small, so that the currently employed CD-R cannot cope with the shortening of the wavelength of the employed laser beam in the future. As a result, it will be highly possible that the information recorded in the currently employed CD-R system cannot be reproduced therefrom by such a future system.\nMany recording materials, for example, the combination of cyanine dye/metal reflection layer, phthalocyanine dye/metal reflection layer, and azo metal chelate dye/metal reflection layer, are proposed for the conventional CD-R. To be more specific, as the recording material for the CD-R, the combination of a cyanine dye/metal reflection layer is proposed in Japanese Laid-Open Patent Applications 1-159842, 2-42652, 2-13656, and 2-168446; a phthalocyanine dye as the recording material is proposed in Japanese Laid-Open Patent Applications 1-176585, 3-215466, 4-113886, 4-226390, 5-1272, 5-171052, 5-116456, 5-69860 and 5-139044; and an azo metal chelate dye as the recording material is proposed in Japanese Laid-Open Patent Applications 4-46186, 4-141489, 4-361088, and 5-279580.\nHowever, there is no write once optical recording medium that can solve the above-mentioned problems."}
-{"text": "Refractory metals such as tantalum and titanium have been used as diffusion barriers, adhesion layers or X-ray lithography masks in electronic applications. Unfortunately, these materials oxidize easily when annealed causing reliability problems due to the increased resistivity and stress. Other refractory metals such as Nb, V, Zr, Hf, Cr, Mo and W also oxidize easily in annealing ambients containing minimal amounts of oxygen and at relatively low temperatures. It is an expensive and difficult process to remove all oxygen from an annealing environment and some processes actually require an oxidizing ambient. Small amounts of oxygen incorporated interstitially in the refractory metals can cause large increases in compressive stress. For instance, in a Ta film that has 10% oxygen incorporation, an increase of 4 to 5 GPa in compressive stress and an increase of 63 .mu..OMEGA.-cm in resistivity have been observed. When used in electronic applications, such increases in resistivity and stress can cause major reliability problems such as poor contact resistance, peeling and cracking.\nOther researchers have utilized refractory metals in electronic applications as diffusion barriers, i.e., Ta diffusion barrier between Si and Cu was used by Holloway, et al., J. Appl. Phys., 71 (11), 5433 (1992). The use of Ta diffusion barrier between a high dielectric material and Si in the fabrication of a high dielectric oxide material has been shown in Grill, et al., J. Mater. Res., 7 (12), 3260 (1992). In the latter work, the fabrication of high dielectric constant oxide materials requires high temperature processing (&gt;650.degree. C.) in oxygen ambients which generally cause the underlying Si to oxidize, and thus creating an additional in-series low dielectric capacitor.\nOther diffusion barrier materials such as TiN, WN and TaN have been used in VLSI applications, i.e., in contact holes for separating conductors from underlying silicon. However, these materials are not suitable oxygen diffusion barriers because they cannot withstand oxidation anneal cycles that the devices may be subjected to. In addition, other researchers have investigated elemental metals such as Pt, Au and Ru for the prevention of diffusion of oxygen to the underlying layer of silicon and its subsequent oxidation. It was found that none of the pure metals prevented the diffusion and the resulting SiO.sub.2 formation. A break in the electrical conduction path to the silicon substrate occurred as a result.\nIn certain semiconductor processes, high dielectric constant Perovskite compounds (such as PZT, PLZT or Ba.sub.x Sr.sub.1-x TiO.sub.3) are deposited onto a substrate. These materials require high temperature (&gt;650.degree. C.) oxygen anneal in order to crystallize. It also requires a Pt seed layer for crystallization into a Perovskite phase which has the highest dielectric constant, i.e., 380 vs. about 40. For simple process integration, it is desirable to use Si as the lower electrode. A problem thus incurred during an oxygen anneal of the Perovskite compound/Pt/Si structure is the formation of an SiO.sub.2 layer at the Pt/Si interface which reduces the effective dielectric constant. A diffusion barrier layer is therefore needed for preventing oxygen diffusion down to the Si layer, and also for preventing the Si diffusion up to the Pt layer.\nIt is therefore an object of the present invention to provide a diffusion barrier for use in a semiconductor device that does not have the shortcomings of the prior art diffusion barriers.\nIt is another object of the present invention to provide a diffusion barrier that can be fabricated in a simple manufacturing process.\nIt is a further object of the present invention to provide a diffusion barrier that is capable of sustaining the high processing temperatures encountered in semiconductor processing steps.\nIt is another further object of the present invention to provide an oxygen diffusion barrier for devices incorporating high dielectric constant oxide layers.\nIt is yet another object of the present invention to provide an oxygen diffusion barrier of a multi-layer aluminum/refractory metal material such that the refractory metal layer does not oxidize (remains electrically conducting) in a subsequent annealing step.\nIt is still another object of the present invention to provide an oxygen diffusion barrier that is a multi-layer thin film of aluminum/refractory metal such that when the film is exposed to atmosphere a thin layer of Al.sub.2 O.sub.3 forms on the surface of aluminum."}
-{"text": "Certain computer operating systems and applications are designed to load modules that are internally requested by the operating system, internally requested by an application, or externally requested by a supported process. Such modules may include, for example, a procedure, a function, a script, a plug-in module, a driver, or a portion of an executable application.\nIn some computer operating systems and/or applications, modules that can be loaded are classified according to a module type. For example, in an operating system, the module type may be used to determine how to properly support the module. In some cases, a sub-system of the operating system, such as a dynamically loadable kernel module (\u201cDLKM\u201d) sub-system, for example, may be used to support modules of different types.\nFIG. 1 illustrates a prior operating system configuration 100 for handling dynamically loadable modules. A table 110 of defined module types and associated type-specific support module identifiers are included in a statically configured operating system kernel 105. Support module identifiers may refer to support modules, such as support module 120, that contain, for example, pre-loading logic, post-loading logic, pre-registration logic, post-registration logic and registration logic.\nWhen the operating system 105 receives a request to load a module 115, the operating system uses the module type definition table 110 to determine whether the module type is supported. Upon locating the correct module type entry in the table, the operating system utilizes the support modules (e.g., support module 120) identified in the table 110 to properly support the requested module 115.\nThe operating system 105 may provide module type-specific support for only the module types defined by the table within the statically configured kernel. This type of operating system configuration is unable to dynamically load modules of types that have not been previously defined in the statically configured portion of the operating system kernel.\nAccordingly, whenever a new module type is to be defined, the statically configured kernel must be reconfigured to include a definition of the new module type. The statically configured kernel must be reloaded to accommodate the newly added module type definition. The reconfiguration of the statically configured kernel results in an increased kernel size and a reloading or rebooting of the operating system."}
-{"text": "1. Field of the Invention\nDisclosed herein is a flashing beacon. More specifically, disclosed herein is a flashing beacon that may be positioned on most any roadway sign or signpost, and that may include a signal unit, a control unit, a solar collector, an activation device (e.g., a timer, microwave emitter, radio transmitter, step-pad, a pushbutton, infrared transmitter, wireless transmitter or like device) and various other accessories.\n2. Reference to Related Art\nAccording to the U.S. Manual on Uniform Traffic Control Devices, flasher mechanisms associated with traffic control signs (e.g., a yield or crosswalk sign) must be positioned on the sign (or signpost) so that flashing signal is about 12 feet above the pavement. The flashing signal must also be programmed or otherwise set to flash continuously at a rate of not less than 50 nor more than 60 times per minute. See MUTCD, Section 4D.11. However, while the guidelines set forth in the uniform regulations are intended to provide a visible warning to drivers, recent testing has suggested that only a small percentage of the public responds to flashing signals that operate according to the uniform regulations. Specifically, recent testing has suggest that only 25% of the public complies with or otherwise responds to flashing signals associated with roadway signage. Therefore, it would be advantageous to have an improved flashing beacon system that may be used with existing or future roadway signage to garner a greater response from the vehicle driving public."}
-{"text": "Certain entities, such as large corporations, government agencies, or universities, may maintain a backbone network (also referred to as a \u201ccomputer network\u201d) to interconnect entity assets. For example, a corporate backbone network may be used to connect different data centers maintained by a corporation. Backbone networks may additionally include one or more points of presence for connecting the backbone network to the Internet. Backbones networks are themselves made up of connectivity devices, such as hubs, switches, and routers, as well as cables to connect the connectivity devices. These components of the backbone network provide the network capacity by which the network demand between assets (e.g., from data center to data center) may be satisfied.\nTo increase backbone network capacity, an entity's network planners may engage in both short-term operational planning and long-term strategic planning. In the short-term, network planners may, for example, acquire additional capacity from an existing cable in the backbone network that has unused capacity, such as by paying for increased cable capacity. In the long-term, network planners may utilize new cables in the backbone network, thereby adding new routes between data centers in the network or augmenting existing routes between data centers.\nTraditionally, the long-term planning done by network planners has been driven by the availability of new supply. That is, as new cables have been laid, network planners have had to decide whether to negotiate for use of the new cables, either through leasing access to the cables or joining a consortium of entities that manage the cable. These decisions have typically been made by network planners based on the current network demands of their backbone networks. In other words, the decisions behind long-term strategic planning have been made by network planners based primarily on the state of the backbone network at that particular time. Such limited input into long-term decision making is inherently inefficient, and risks the ability of the backbone network to meet future network demands. The currently available tools for network planning are inefficient. The planning of network can include planning a transport layer of the network, e.g., physical link layer such as laying of fiber optics, their capacities, and an Internet Protocol (IP) layer, e.g., placement of routers in the network, how the routers have to be connected. The current tools typically facilitate planning one layer or the other and do not support multi-layer planning. It would therefore be beneficial to facilitate strategic planning of a backbone network."}
-{"text": "This specification relates to scoring images that are related to particular entities.\nThe Internet provides access to a wide variety of resources, for example, webpages, images, audio files, and videos. A variety of search engines are available for identifying particular resources accessible over the Internet. For example, digital images that satisfy a user's informational need can be identified by an image search process in which keywords or other data are processed to identify collections of digital images. Each image is identified to the user by an image search result. The image search result may include a thumbnail of an image to which it links.\nSearch results generated by a search engine are generally ranked and presented to a user in an ordered list of search results. The rankings may be based on how the search engine determines the relevance of resources specified by the search results to a particular query issued by the user.\nA very common entity that people search for is a person. Thus, users use search engines to identify images associated with a particular person. The images associated with the particular person typically are images that depict the face of the particular person."}
-{"text": "1. Technical Field of the Invention\nThe present invention relates to the decoding of radiofrequency transmission channels conveying coded digital information.\nThe invention thus applies advantageously to digital terrestrial television, that is to say using signals transmitted between television antennas, as defined, for example, in the European DVB-T (Digital Video Broadcasting-terrestrial) or in the North-American ATSC (Advisory Committee on Advanced Television Service Digital TV) specification, or to digital cable television, as defined, for example, in the European DVB-C (Digital Video Broadcasting-cable) or in the North-American ITU-J83 Annex A/B/C/D (International Telecommunication Union) specification, all these telebroadcasts being based on the MPEG transmission standards, and use for example to convey information, quadrature digital modulation, or else modulation of the COFDM type according to terminology well known to the person skilled in the art.\nThe invention also relates in particular to tuners.\n2. Description of Related Art\nThe television signals received at the input of a tuner of a receiver are composed of the entire set of channels transmitted that lie in the 45 MHz-860 MHz frequency band. The objective of the tuner is to select the desired channel and to output a baseband signal on the in-phase path (I path) and on the quadrature path (Q path). This signal is thereafter converted into a digital signal and demodulated. The channel decoding processing thereafter also comprises a block which distinguishes, typically by means of majority logic, zeros from ones, then performs all the error correction, that is to say typically a Viterbi decoding, the interleaving, Reed-Solomon decoding and deshuffling. The channel decoding device outputs packets that are decoded in a conventional manner in a source decoding device in accordance with the MPEG standards so as to reoutput the initial audio and video signals transmitted via antennas or via the cable.\nCurrent studies are aimed at researching ever more integrated solutions for the embodiment of digital terrestrial or cable television receivers. However, this research hits technological difficulties related to the embodiment of the means making it possible, on the one hand, to correctly discern a channel from the entire set of channels present, and, on the other hand, to correctly discriminate the signal from noise. Specifically, whereas in digital satellite television, all the channels have nearly the same power, this is not the case in particular in digital terrestrial television. Thus, it is for example possible to have adjacent channels exhibiting a very high power with respect to the desired channel, for example a difference in power of the order of 40 dB. It is therefore necessary to be able to ensure rejection of 40 dB on the adjacent channels. Moreover, the specification imposes a signal/noise ratio of the order of 30 dB. Hence, in the worst case, it is therefore essential for it to be possible to bring the level of an adjacent channel to 70 dB below its input level. This results in extremely significant constraints on the filters, this being extremely constraining in respect of integrated solutions.\nThere exists a need to solve this problem. More specifically, there is a need for a totally integrated solution for the tuner, given the characteristics of digital television signals, in particular terrestrial ones, exhibiting more significant rejection constraints than digital cable television signals."}
-{"text": "1. Field of the Invention\nThe present invention relates to an adjusting apparatus for automatically adjusting positional deviation of stereoscopic images by means of geometrical image transformation processes.\n2. Discussion of the Background Art\nIn recent years, a stereoscopic vehicle surrounding monitoring apparatus using a pair of left and right cameras (stereoscopic camera having solid image element like CCD) mounted on the vehicle has been interested by automobile engineers. To detect a distance to an object, first respective picture element blocks having coincidence of brightness are found in left and right images (stereo matching), then distance data are calculated according to the principle of triangulation from a parallax, namely a relative deviation amount, between both picture element blocks. Consequently, in order to calculate distance data with high reliability, it is desirable that there exists no positional deviation other than the parallax in a pair of left and right images (stereoscopic images). In actual world, however, the stereoscopic camera has some amount of positional error caused when the camera is installed on the vehicle.\nWith respect to this, Japanese Patent Application Laid-open No. Toku-Kai-Hei 10-307352 discloses a technology in which the positional deviation of the stereoscopic camera is corrected by applying a geometric transformation to the stereoscopic image. That is, when an initial adjustment of the positional deviation is made or when a readustment of the positional deviation generated by aged deterioration is made, a dedicated correction detecting device is connected with an image correction apparatus performing the affine transformation to calculate the difference of angle of view, a rotational deviation or a parallel deviation of the stereoscopic image obtained by imaging a specified pattern for adjustment and to establish (reestablish) parameters of the affine transfomation according to the result of the calculation. The positional deviation is equivalently corrected by applying the affine transformation to images based on thus established affine parameters.\nHowever, according to the prior art, a special adjustment pattern is imaged by the stereoscopic camera and the deviation is corrected based on the position of the pattern in images. Accordingly, when the correction is performed, it is necessary to interrupt the ordinary surroundings monitoring control and as a result this prior art is not suitable for a real time processing in which the monitoring control is carried out concurrently.\nIt is an object of the present invention to provide an adjusting apparatus capable of correcting the positional deviation of stereoscopic images in real time in parallel with the surroundings monitoring control.\nIn order to attain the object, a positional deviation adjusting apparatus of a stereoscopic camera taking a first image or a reference image and a second image or a comparison image, comprises an image correcting means for transforming an image geometrically according to transformation parameters, a reference area establishing means for establishing a plurality of reference areas in a first image transformed by the image correcting means, a reference point establishing means for establishing a reference point in the reference area for each of the reference points, a reference line establishing means for establishing a reference line passing through the reference points, a searching area establishing means for establishing a searching area corresponding to the reference area in a second image, a correlation point identifying means for identifying a correlation point having a correlation with the reference point for each of the searching areas, an approximation line calculating means for calculating an approximation line based on the correlation points and a parameter correcting means for correcting the transformation parameters such that the approximation line coincides with the reference line."}
-{"text": "Hydrocracking is a versatile petroleum refining process, which enjoys widespread and growing use in the refining industry. Hydrocracking has the ability to process a wide range of difficult feedstocks into a variety of desirable products. Feedstocks, which may be treated by this process, include heavy naphthas, kerosenes, refractory catalytically cracked cycle stocks and high boiling virgin and coker gas oils. At high severities, hydrocracking can convert these materials to gasoline and lower boiling paraffins; lesser severities permit the higher boiling feedstocks to be converted into lighter distillates such as diesel fuels and aviation kerosenes.\nHydrocracking is conventionally carried out at moderate temperatures of 350\u00b0 C. to 450\u00b0 C. (650\u00b0 F. to 850\u00b0 F.) and at high pressures, over 7,000 kPa (1,000 psig) because the thermodynamics of the hydrocracking process become unfavorable at higher temperatures. In addition, high hydrogen pressures, usually at least 800 psig are usually required to prevent catalyst aging and so to maintain sufficient activity to enable the process to be operated with a fixed bed of catalyst for periods of one to two years without the need for regeneration.\nEnvironmental concerns, especially with sulfur oxides and nitrogen oxides emissions, have led petroleum refiners to depend more heavily than in the past on hydrodesulfurization and hydrocracking processes. Availability of by-product hydrogen from naphtha reforming no doubt has also cooperated to foster this dependence. Other factors too, have come into play to make hydroprocessing of increasing importance. Among these factors is that high quality crude oils for lube and fuels refineries are expected to progressively become more scarce. Also, refineries that include a fluid catalytic cracking (FCC) plant generate large volumes of dealkylated, aromatic refractory effluents, commonly known as FCC Cycle Oils. Decrease in demand for the fuel oil products into which these FCC Cycle Oils were previously incorporated has added to the practice of working them off by incorporation with a hydrocracker feedstock. The hydrocracking process, unlike catalytic cracking, is able to effectively upgrade these otherwise refractory materials.\nHydrocracking is an established petroleum refining process. The hydrocracking feedstock is invariably hydrotreated before being passed to the hydrocracker in order to remove sulfur and nitrogen compounds as well as metals and, in addition, to saturate olefins and to effect a partial saturation of aromatics. The sulfur, nitrogen and oxygen compounds may be removed as inorganic sulfur, nitrogen and water prior to hydrocracking although interstage separation may be omitted, as in the Unicracking-JHC process. Although the presence of large quantities of ammonia may result in a suppression of cracking activity in the subsequent hydrocracking step, this may be offset by an increase in the severity of the hydrocracking operation.\nIn the hydrotreater, a number of different hydrogenation reactions take place including olefin and aromatic ring saturation but the severity of the operation is limited so as to minimize cracking. The hydrotreated feed is then passed to the hydrocracker in which various cracking and hydrogenation reactions occur.\nIn the hydrocracker, the cracking reactions provide olefins for hydrogenation while hydrogenation in turn provides heat for cracking since the hydrogenation reactions are exothermic while the cracking reactions are endothermic; the reaction generally proceeds with generation of excessive heat because the amount of heat released by the exothermic hydrogenation reactions usually is much greater than the amount of heat consumed by the endothermic cracking reactions. This surplus of heat causes the reactor temperature to increase and accelerate the reaction rate, but control is provided by the use of hydrogen quench.\nConventional hydrocracking catalysts combine an acidic function and a hydrogenation function. The acidic function in the catalyst is provided by a porous solid carrier such as alumina, silica-alumina, or by a composite of a crystalline zeolite such as faujasite, Zeolite X, Zeolite Y or mordenite with an amorphous carrier such as silica-alumina. The use of a porous solid with a relatively large pore size in excess of 7 A is generally required because the bulky, polycyclic aromatic compounds which constitute a large portion of the typical feedstock require pore sizes of this magnitude in order to gain access to the internal pore structure of the catalyst where the bulk of the cracking reactions take place.\nThe hydrogenation function in the hydrocracking catalyst is provided by a transition metal or combination of metals. Noble metals of Group VIIIA of the Periodic Table, especially platinum or palladium may be used, but generally, base metals of Groups IVA, VIA and VIIIA are preferred because of their lower cost and relatively greater resistance to the effects of poisoning by contaminants. The preferred base metals for use as hydrogenation components are chromium, molybdenum, tungsten, cobalt and nickel; and, combinations of metals such as nickel-molybdenum, cobalt-molybdenum, cobalt-nickel, nickel-tungsten, cobalt-nickel-molybdenum and nickel-tungsten-titanium have been shown to be very effective and useful.\nThe porous support, which provides the acidic functionality in the catalyst, may comprise either an amorphous or a crystalline material or both. Amorphous materials have significant advantages for processing very high boiling feeds which contain significant quantities of bulky polycyclic materials (aromatics as well as polynaphthenes) since the amorphous materials usually possesses pores extending over a wide range of sizes and the larger pores, frequently in the size range of 100 to 400 Angstroms are large enough to provide entry of the bulky components of the feed into the interior structure of the material where the acid-catalyzed reactions may take place. Typical amorphous materials of this kind include alumina and silica-alumina and mixtures of the two, possibly modified with other inorganic oxides such as silica, magnesia or titania.\nCrystalline materials, especially the large pore size zeolites such as zeolites X and Y, have been found to be useful for a number of hydrocracking applications since they have the advantage, as compared to the amorphous materials, of possessing a greater degree of activity, which enables the hydrocracking to be carried out at lower temperatures at which the accompanying hydrogenation reactions are thermodynamically favored. In addition, the crystalline catalysts tend to be more stable in operation than the amorphous materials such as alumina. The crystalline materials may, however, not be suitable for all applications since even the largest pore sizes in these materials, typically about 7.4 Angstroms in the X and Y zeolites, are too small to permit access by various bulky species in the feed. For this reason, hydrocracking of residuals fractions and high boiling feeds has generally required an amorphous catalyst of rather lower activity.\nThe crystalline hydrocracking catalysts generally tend to produce significant quantities of gasoline boiling range materials (approximately 330\u00b0 F.-, 165\u00b0 C.-) materials as product. Since hydrocracked gasolines tend to be of relatively low octane and require further treatment as by reforming before the product can be blended into the refinery gasoline pool, hydrocracking is usually not an attractive route for the production of gasoline. On the other hand, it is favorable to the production of distillate fractions, especially jet fuels, heating oils and diesel fuels since the hydrocracking process reduces the heteroatom impurities characteristically present in these fractions to the low level desirable for these products.\nU.S. Pat. No. 8,372,772, issued Feb. 12, 2013, discloses a hydrocracking catalyst comprising zeolite crystallized as a layer on the surface of a porous alumina-containing matrix, the matrix being shaped by extrusion and arranged in a configuration to provide macropores in which the zeolite layer is provided on the walls of the macropores. The in-situ produced Y-fauajasite on matrix provides acidic functionality in the catalyst and is also used as a support for hydrogenating metals.\nU.S. patent applications (Pub. Nos. 2004/0220046 A1 and 2004/0235642 A1) teach preparation methods for structurally enhanced zeolite microsphere FCC catalysts where zeolite crystallizes on the walls of the pores of the matrix.\nAs is known in the industry, the hydrocracking process is a fixed-bed process with catalysts in shaped forms of extrudates, pellets or spheres. The particle size is much larger than FCC microspheres. While the above-mentioned U.S. Pat. No. 8,372,772 discloses in-situ crystallization of pre-shaped particles, such as extrudates, the patent does not otherwise provide a specific method which provides sufficient exchange of nutrients between the liquid phase and the solid phase to provide crystallization and allows the zeolite to form on the pore walls, as well as maintain the porosity and integrity of the extrudate."}
-{"text": "The present invention relates to photopolymer materials sensitive to infrared, near infrared, red and green light radiation for initiating polymerization and to applications of such photopolymer, like holographic polymer dispersed liquid crystal (H-PDLC) or reversible dye doped photopolymer (RDDP) materials, for making optical devices. The invention relates to holographic polymer dispersed liquid crystal and reversible dye materials having improved electrical and optical switching properties.\nInfrared (IR) diode sources are largely used in integrated photonic circuits and the list of their applications grows very rapidly. The management of the radiation of these sources requires the fabrication of optical elements such as lenses, interconnects, modulators, etc. Holographic diffractive elements are emerging as very promising for these applications. These include wavelength selective holographic interconnects (1 to N, or M to N), couplers, lenses, mirrors. Today, a new class of Holographic Photopolymer Dispersed Liquid Crystal Materials (H-PDLCs) is considered to be one of the most viable technologies for the development of reflective color displays, switchable holographic optical elements (such as Bragg gratings for Wavelength Division Multiplexing (WDM) devices), switchable-focus lenses, etc. See for example Crawford et al., xe2x80x9cReflective color LCDs based on H-PDLC and PSCT technologiesxe2x80x9d, J Soc. Information Display, 1996, 5(1); Domash et al., xe2x80x9cElectronically switchable waveguide Bragg gratings for WDM routingxe2x80x9d, 1997 Digest of the IEEE. /LEOS Summer Topical Meetings: Vertical-Cavity Lasers/Technologies for a Global Information Infrastructure/WDM Components Technology; and Domash et al., xe2x80x9cSwitchable-focus lenses in holographic polymer dispersed liquid crystalxe2x80x9d, Proceedings of the SPIExe2x80x94The International Society for Optical Engineering, vol. 2689, (Diffractive and Holographic Optics Technology III, San Jose, Calif., USA, 1-2 Feb. 1996.) SPIE-Int. Soc. Opt. Eng, 1996. p. 188-94.\nCommercially available holographic materials are generally sensitive in the UV/visible part of the spectrum only. Thus, the actual fabrication of above-mentioned elements requires an initial recording step, where a UV-Visible laser source is used, and then, a further adaptation or adjustment of the obtained element for utilization with near IR wavelengths (e.g. 800-850 nm or 1300-1500 nm), which are used in local or long distance communication systems. Due to strong astigmatism and divergence of the used diode lasers, this work is difficult and has poor efficiency. It would therefore be highly desirable to provide in situ recording of holographic diffractive elements with lasers, which are already integrated in the given photonic circuit, thus providing self-alignment of the photonic circuits. Thus, there is a need to extend the sensitivity of these materials up to communication wavelengths. Namely, for in situ holographic recording of optical components with diode lasers operating in the 800-855 nm, among which there are Vertical Cavity Surface Emitting Lasers (VCSEL), it is important to have H-PDLC material with suitable holographic characteristics, which include high sensitivity and diffraction efficiency of the recording, low scattering and noise level, and low switching voltage.\nUp to now, the sensitivity of the existing H-PDLC materials has been extended up to 790 nm (Natarajan et al., Report on Photonics West 98, San Jose, January 1998). Photopolymerizable materials have been recently proposed with sensitivity in the 800-850 nm area, for example in co-pending application Ser. No. 09/503,207, filed on Feb. 14, 2000. See also EP 0 223 587; EP 0 387 087; EP 0 389 067; U.S. Pat. No. 4,343,891; Chatterjee et al., J. Am. Chem. Soc., 1990, 112, 6329; Schuster et l, Photochem. Photobiol. A: Chem., 1992, 65, 191; Chatterjee et al., J. Am. Chem. Soc., 1988, 110, 2326; Cooper et al., J. Am. Chem. Soc., 1963, 85, 1590; and Noiret et al., Pure and Applied Optics, 1994, 3(1), 55-71. Imaging applications with low resolution (such as printing plates) were also successfully explored. Some of the materials were the subject of the study for holographic gratings recording and only very low level of performance was achieved (7% of diffraction efficiency at sensitivity of about 300-500 mJ/cm2), as reported in the mentioned Noiret et al. article supra. This low level of performance makes such polymers impractical for commercial application. All the above materials still suffer from limitations such as spatial resolution, diffraction efficiency etc.\nIn accordance with the present invention, there is now provided a photopolymerizable formulation sensitive to light in a green to infrared region of the optical spectrum comprising:\na) a photopolymerizable monomer, oligomer or mixtures thereof;\nb) a photoinitiator sensitive to light in the green to infrared region;\nc) an additive for increasing the refractive index and decreasing the viscosity of the formulation; and\nd) an optional filler having optical properties selected to contrast with optical properties of a polymer resulting from photopolymerization of the monomer, oligomer or mixtures thereof.\nPreferably, the formulation comprises filler, which is a liquid crystal or a reversible dye, or combinations thereof, or any other filler materials such as mesogens having polar or functional groups. The liquid crystal preferably has a polar group attached to ends of the molecule chain elements, whereby droplets of LC matter are more efficiently formed in the polymer dispersed in the liquid crystal material. Larger droplets allow for more efficient optical state switching.\nIn a second aspect of the invention, there is provided a process for making an optical device comprising the steps of:\npreparing a photopolymerizable formulation sensitive to light from green to infrared region of the optical spectrum, the formulation comprising\na) a photopolymerizable monomer, oligomer or mixtures thereof;\nb) a photoinitiator sensitive to light in the region\nc) an additive for increasing the refractive index and decreasing the viscosity of the formulation; and\nd) an optional filler having optical properties selected to contrast with optical properties of a polymer resulting from photopolymerization of the monomer, oligomer or mixtures thereof;\napplying a layer of the formulation on an optical element;\nfirst exposing the optical element comprising the formulation to a light source emitting light in the region, to polymerize the formulation and produce a recording pattern on the optical element; and\noptionally second exposing the optical element to the light to polymerize any remaining portion of the formulation not polymerized during first exposing, thereby producing said optical device.\nFinally, in a third aspect of the present invention, there is provided a process of making an optical connection between at least two waveguides or light guides of an optical element by photopolymerization, the process comprising the steps of:\npreparing a photopolymerizable formulation sensitive to light from green to infrared region of the optical spectrum, the formulation comprising:\na) a photopolymerizable monomer, oligomer or mixtures thereof;\nb) a photoinitiator sensitive to light in the green to infrared region;\nc) an additive to increase the refractive index and decrease the viscosity of the formulation; and\nd) an optional filler;\napplying a sufficient amount of the formulation between the waveguides or light guides of the optical element to be connected; and\nexposing the optical element to a light source emitting light in the region to polymerize the formulation, thereby making an optical connection between the at least two waveguides or light guides.\nPreferably, the light is transmitted during the exposing through at least one of the at least two waveguides or light guides. While a filler is also optional in the case of optically coupling waveguides or light guides, a filler having the desired optical properties is preferred to reduce the amount of polymer used and to impart desirable physical properties to the formulation."}
-{"text": "This invention relates to electronic circuitry and, more particularly, to a charge pump for producing a negative substrate bias in a complementary metal oxide semiconductor (CMOS) integrated circuit.\nMOS transistors are commonly used in electronic circuits such as dynamic random access memories (DRAMS). In an NMOS transistor, an N-type source region is separated from an N-type drain region by a P-type channel region. All three regions are formed in a P-type semiconductor substrate. By applying a positive voltage to a gate electrode disposed above the channel region, electrons gather in the channel region between the source region and the drain region to allow current to flow from the drain region to the source region. PMOS transistors have the same structure except the conductivity types of the various regions are reversed and a negative gate voltage is required to allow current to flow from the source region to the drain region.\nIt has been found that NMOS transistors operate better when the P-type substrate of the NMOS (or of the NMOS transistors in a CMOS circuit) is driven negative with respect to circuit ground, in other words there is a negative substrate bias. Such a negative substrate bias provides a number of advantages in terms of the overall circuit performance. More specifically, a negative substrate bias decreases the NMOS transistor source and drain capacitance, decreases the likelihood of latchup, decreases PN diode injection when a node is driven below ground, and decreases the effective body effect, all of which are desirable in CMOS circuits.\nTypically a charge pump circuit is used to create the negative substrate bias. Once a negative substrate bias is achieved, however, it does not last forever. For example, when an NMOS transistor is conductive with a relatively high drain to source voltage, some of the electrons traveling from the source region to the drain region collide with atoms in the channel region with enough energy to cause electron/hole pairs to form. The positive gate voltage attracts the generated electrons to the surface of the channel while the positive drain voltage attracts them to the drain where they simply add to the normal flow of electrons from source to drain. The positively charged holes, by contrast, are repelled by the positively charged gate away from the channel region into the substrate. The substrate current created by the excess holes makes the substrate more positively charged, thus counteracting the negative substrate bias. In DRAMS, a substantial amount of substrate current is generated whenever the memory is read or written, since many transistors are switched on and off at that time. This component of substrate current may be orders of magnitude above the background (i.e., standby) leakage current of all the reverse biased P--N diodes throughout the circuit. Therefore, the charge pump must remove low substrate current during standby and high substrate current during high activity to maintain the negative substrate bias.\nFIG. 1 is a conceptual schematic diagram of a charge pump 2 which includes a first switch 4 coupled between a positive power supply voltage (V.sub.cc) and, a first terminal 6 of a capacitance C1. A second switch 8 is coupled between a ground potential (V.sub.ss) and a second terminal 10 of capacitance C1. A third switch 12 is coupled between (V.sub.ss) and terminal 6 of capacitance C1, and a fourth switch 14 is coupled between the substrate (represented by the voltage (V.sub.bb)) and terminal 10 of capacitance C1. In operation, switches 4 and 8 are both closed (made conductive) for charging capacitance C1 to a voltage equal to the difference between (V.sub.cc) and (V.sub.ss). In FIG. 1, (V.sub.cc)=+5 volts and (V.sub.ss)=0 volts, so capacitance C1 charges with node 6 five volts more positive than node 10. Thereafter, switches 4 and 8 are opened and switches 12 and 14 are both closed. Since the positive terminal 6 of capacitance C1 is now coupled to a ground potential, the negative terminal 10 of capacitance C1 tries to drive V.sub.bb to negative 5 volts through switch 14. Thereafter, switches 12 and 14 are opened, and the sequence repeats itself. An oscillator (not shown) typically controls the repetitive switching sequence, and a detector (not shown) monitors the substrate voltage and controls the pumping operation to maintain the substrate at the proper negative voltage level.\nAs discussed in more detail below, known charge pumps consume a substantial amount of power (often 1 milliwatt or more even when no further pumping is required), often work against themselves by adding positive substrate current as they operate, and generally operate inefficiently."}
-{"text": "This invention relates in general to die-casting machines and in particular to high-speed die-casting machines for the production of precision miniature castings where the weight of the produced casting generally does not exceed two ounces and the face of the die generally does not exceed approximately ten square inches. While these size limitations provide a suitable focus, many of the features of the present invention are applicable to much larger machines.\nHeretofore, machines for the die casting of precision miniature pieces have been generally of the Dynacast type, a die-casting machine developed in the mid-1960's and the subject of several United States patents, including U.S. Pat. No. 3,364,981. This machine, at the time of its development, drastically changed the manufacture of miniature die-casting pieces which prior to the time of the U.S. Pat. No. 3,364,981 were manufactured in large conventional die-casting machines using multiple-cavity dies. However, the state of the art has remained practically unchanged during the last 20 years. At the expiration of the above-mentioned patent, the machine has been duplicated by several manufacturers with more or less the same features but without eliminating the various shortcomings.\nOne shortcoming of the Dynacast machine is the fact that the machine can use only two vertical cores and only occasionally a frontal core. The zinc injection system of the Dynacast machine is made by a gooseneck partially submerged into the liquid zinc of the melting pot. The gooseneck has a swiveling movement in order to bring the nozzle in contact with the die face and to move the nozzle away from the die face to break the runner from the liquid metal. Besides the mechanical elements needed for this movement (which adds complexity to the machine), this system reduces the speed at which the machine can operate. The position of the gooseneck precludes the use of a fourth core.\nAnother shortcoming of the Dynacast machine includes the way the level of the zinc in the neck of the gooseneck is controlled. This system leaves a certain amount of air ahead of the zinc in the neck of the gooseneck, which amount of air is variable and depends upon several parameters which are not easily controlled, such as the temperature of dies, size of the shots, speed of the machine and temperature of the nozzle. Another drawback of the Dynacast machine is the way the toggle mechanism operates in order to close the die halves and to create the force between them. This force generated between the dies varies, depending on the wear of the mechanical components of the toggles, temperature of the dies, temperature of the plate of the machine, and so forth.\nA further drawback of the Dynacast machine is the fact that it cannot produce parts with what is known in the industry as a \"hardware\" finish; in other words, parts that do not need to be plated. The reason for such inability involves the way the zinc is injected into the die, the gating (which has to be small because of the subsequent separation from the part with the tumbling operation) and the system of heat transfer on the die. The parts from the Dynacast machine are ejected with the gates attached to them, and this requires a subsequent operation (tumbling) in order to separate the parts from the gate, with all of the negative attributes of such operation including surface pitting, dimensional alteration, and breaking of sharp corners.\nThe present invention provides a machine which is designed to overcome the above shortcomings and disadavantages. In order to achieve this object, the disclosed machine produces castings which may require up to four cores along the plane of the dies (substantially two vertical and two horizontal) and two cores in a plane perpendicular to the surface of the die with a total of six cores. The machine of the present invention produces castings with much less porosity than any other machine used today and with surface finish of the type as previously mentioned which is known in the industry as a hardware finish. The cast pieces can be trimmed automatically in the machine so that the pieces are separated from the gate, the gates being discharged in a separate container for a subsequent remelt.\nAnother reference of possible interest to the heat transfer aspects of the present invention is U.S. Pat. No. 4,248,289 which issued to Perella et al. This patent discloses a means to heat the heat-exchange medium, water, by electric immersion heaters to preheat the die aove the normal boiling point of water by controlling the internal pressure of the die-cooling cavity. Die cooling is accomplished by the heat of evaporation at the temperature and pressure required for best casting quality. The wall thickness of the die cavity to water cavity is sized according to the heat-transfer rate or the die-cavity material. As steam is produced within the cavity, make-up water is introduced to maintain a wetted surface.\nIn die-casting practice, the \"coring\" of a cast part encompasses many forms. The core may be in the form of an entry hole, round, square or otherwise, a hollow space, a countersunk or extended boss, an embossed face, or any shape that is required that is not formed by the two-part mold cavity. Some core requirements are parallel to the line of separation of the mold halves but are independent in movement. Other cores not in the line of mold separation movement must be positioned and supported either by devices integral with the die blocks or as a function of the die-casting machine. Those arrangements which are integral with the molding dies add substantially to the cost of the dies and are not practical for miniature casting work. Since the hydraulic pressure area of some cores can be equal to that of the die cavity, it is necessary to have equal support to that of the die cavity to lock the core in a fixed position. This may be accomplished by toggle movement to rapidly position the core in the die and then lock the position as securely as the main die-closing action."}
-{"text": "Ovarian cancer is the eighth most common type of cancer among women. The American Cancer Society estimates that about 22,220 new cases of ovarian cancer will be diagnosed in the United States during 2005. Ovarian cancer accounts for about 3% of all cancers in women. Because many ovarian cancers cannot be detected early in their development, they account for a disproportionate number of fatal cancers, being responsible for almost half the deaths from cancer of the female genital tract; more deaths than any other reproductive organ cancer. Older women are at higher risk. More than half of the deaths from ovarian cancer occur amongst women between 55 and 74 years of age. About 25% of ovarian cancer deaths occur amongst women between 35 and 54 years of age.\nThe main treatments for ovarian cancer are surgery, chemotherapy, and radiation therapy. Combinations of these treatments are used to treat ovarian cancer.\nSurgery is the usual initial treatment for women diagnosed with ovarian cancer. The ovaries, the fallopian tubes, the uterus, and the cervix are usually removed. Staging during surgery (to find out whether the cancer has spread) generally involves removing lymph nodes, samples of tissue from the diaphragm and other organs in the abdomen, and fluid from the abdomen. If the cancer has spread, the surgeon usually removes as much of the cancer as possible. This reduces the amount of cancer that will have to be treated later with chemotherapy or radiation therapy\nChemotherapy is the use of drugs to kill cancer cells. Chemotherapy may be given to destroy any cancerous cells that may remain in the body after surgery, to control tumour growth, or to relieve symptoms of the disease. Most drugs used to treat ovarian cancer are given intravenously or directly into the abdomen through a catheter.\nRadiation therapy, also called radiotherapy, involves the use of high-energy rays to kill cancer cells. Radiation therapy affects the cancer cells only in the treated area. The radiation may come from a machine or women receive a treatment called intraperitoneal therapy in which radioactive is put directly into the abdomen through a catheter.\nDeciding on a particular course of treatment is typically based on a variety of prognostic parameters and markers [Fitzgibbons, et al. (2000) Arch. Pathol. Lab. Med. 124:966-978; Hamilton and Piccart (2000) Ann. Oncol. 11:647-663)], including genetic predisposition markers BRCA-1 and BRCA-2 [Robson (2000) J. Clin. Oncol. 18:113sup-118sup].\nAlthough many ovarian cancer patients are effectively treated, the current therapies can all induce serious side effects, which diminish quality of life. Moreover, approximately 85% of the patients that have been effectively treated with platinum- and paclitaxel-based chemotherapy, including complete responses, relapse within two years after treatment.\nThe identification of novel therapeutic targets is essential for improving the current treatment of ovarian cancer patients. Recent advances in molecular medicine have increased the interest in tumour-specific cell surface antigens that could serve as targets for various immunotherapeutic or small molecule strategies.\nAmong the various elements of the immune system, T lymphocytes are probably the most adept to recognize and eliminate cells expressing foreign or tumour-associated antigens. Cytotoxic T Lymphocytes (CTLs) express the CD8 cell surface marker and are specialized at inducing lysis of the target cells with which they react via the perforin/granzyme and/or the Fas/Fas-L pathways. The T-cell receptor (TCR) for antigen of CTLs binds to a molecular complex on the surface of the target cell formed by small peptides (8-11) residues derived from processed foreign or tumour associated antigens, which associate with major histocompatibility complex (MHC) class I molecules.\nThe other major T-cell subset, helper T lymphocytes (HTLs or T helper cells), is characterized by the expression of CD4 surface marker. The T helper cells recognize slightly larger peptides (11-20 residues), also derived from foreign or tumour associated antigens, but in the context of MHC class II molecules, which are only expressed by specialized antigen presenting cells (APCs) such as B lymphocytes, macrophages and dendritic cells (DCs).\nAs a consequence of TCR stimulation of naive CTLs and HTLs by peptide/MHC complexes on APCs, the CTLs mature into effector killer cells capable of lysing (tumour) cells that express the corresponding peptide/MHC class I complex. Activated HTLs amplify CTL responses by making the APCs more effective at stimulating the naive CTLs and by producing lymphokines that stimulate the maturation and proliferation of CTLs. The potentiating effect of T helper cells occurs both in secondary lymphoid organs where the immune response is initiated and at the tumor site where CTL responses need to be sustained until the tumour cells are eliminated. Thus, one would predict that vaccines should stimulate both tumour-reactive CTLs and HTLs to generate effective antitumour immunity.\nAntigens suitable for immunotherapeutic cancer strategies should be highly expressed in cancer tissues and ideally not in normal adult tissues. Expression in tissues that are dispensable for life, however, may be acceptable.\nA number of antigens suitable for immunotherapeutic strategies in the treatment of ovarian tumors have been described so far, including MUC1, CTs, SP17 and Her2/neu.\nPolymorphic epithelial mucin (MUC1) is a transmembrane protein, present at the apical surface of glandular epithelial cells. It is often overexpressed in ovarian cancer (in more than 90% of all ovarian cancers), and typically exhibits an altered glycosylation pattern, resulting in an antigenically distinct molecule. MUC1 is in early clinical trials as a vaccine target [Gilewski, et al. (2000) Clin. Cancer Res. 6:1693-1701; Scholl, et al. (2000) J. Immunother. 23:570-580]. The tumour-expressed protein is often detectable as tumor marker in the circulation [cf. Bon, et al. (1997) Clin. Chem. 43:585-593].\nA unique class of differentiation antigens, the cancer/testis (CT) antigens, are not expressed in normal tissues except for testis and, in some cases, placenta. This fact makes CT antigens attractive targets for specific immunotherapy of cancer. The function of the majority of the CT antigens is currently unknown. Tammela et al. [Tamella, et al. (2004) Cancer Immunity 4:10-21] demonstrated that SCP-1, a CT antigen with a known role in gamete development, is expressed in 15% of ovarian cancer cases. It was suggested that because of its restricted expression in normal tissues and its aberrant expression in tumour tissues SCP-1 might serve as a potential target for vaccine therapy in ovarian cancer.\nAnother potential target for immunotherapy in patients with ovarian carcinoma is the sperm protein 17 (SP17). Sp17 was found to be expressed in the primary tumor cells from 70% of the patients with ovarian carcinoma. The restricted expression of Sp17 in normal tissue makes it an ideal target for tumour vaccine. A recombinant Sp17 protein was used with monocyte-derived dendritic cells and autologous peripheral blood mononuclear cells to generate Sp17 specific cytotoxic T-lymphocytes (CTLs). Human leukocyte antigen (HLA) class II restricted Sp17 specific CTLs were generated successfully from the peripheral blood of three patients with ovarian carcinoma at the time of disease presentation. These CTLs were able to lyse autologous Epstein-Barr virus-transformed lymphoblastoid cells in an Sp17-dependent manner. The CTLs also lysed Sp17-positive autologous tumour cells, suggesting that Sp17 is processed and presented in association with the HLA class I molecules in Sp17-positive tumour cells. [Chiriva et al. (2002) Cancer 94(9):2447-2453]\nHuman epidermal growth factor receptor 2 (Her2/neu) is an oncogene that is activated by gene amplification with the increased expression of another (normal) gene product. Her2/neu is overexpressed in 20 to 30% of patients with breast and ovarian cancer. Initial studies to develop a peptide based HER-2/neu vaccine were performed in a rat model [Disis et al. (1999) Clinical Cancer Research 5:1289-1297]. No T-cell responses or anti-body responses were observed in animals immunized with intact rat neu protein. By marked contrast, tolerance to rat neu protein in rats, could be circumvented by immunization with a peptide based vaccine. Rats immunized with neu peptides designed for eliciting CD4+ T-cell responses, generated T-cell and antibody responses specific for both the immunizing peptides and the whole protein.\nBrossart et al. demonstrated that patients with advanced breast and ovarian cancer could be efficiently vaccinated with autologues dendritic cells (DCs) pulsed with Her2/neu- or MUC1-derived peptides. In 5 out of 10 patients peptide specific CD8+ cytotoxic T lymphocytes could be detected in the peripheral blood. It was reported that MAGE-3- and CEA-peptide-specific CD8+ T cells were observed in one patient treated with MUC-1 peptide-pulsed DCs, and MUC-1 specific T-cells were observed in another patient after vaccination with HER2/neu derived peptides. It was suggested by Brossart et al. [Brossart et al. (2002) Transfus Apher Sci. 27(2):183-186] that this indicated that epitope spreading occurred in these patients upon treatment.\nEpitope spreading is a recognized phenomenon of autoimmune responses and is believed to be an exacerbating factor in CD4+ T cell-mediated autoimmune diseases. The phenomenon has been demonstrated in murine relapsing-remitting experimental autoimmune encephaleomyelitis (EAE), Theiler's murine encephalomyelitis virus-induced demyelineating disease and diabetes in the non-obese diabetic (NOD) mouse. A model has been suggested for how epitope spreading in autoimmune diseases mediated by CD4+ T cells occurs. This model is supported by direct evidence that tissue damage, TCR ligation on CD4+ T cells by MHC class II-peptide complexes, CD40-CD40 ligand interactions and CD28-mediated co-stimulation are required for epitope spreading to become manifest. It is thought that an initiating self-antigen or a persistent viral epitope, presented in MHC class II molecules on the surface of professional antigen-presenting cells (APC) residing in the target tissue, causes the activation of CD4+ T cells specific for that antigen. This T cell activation results in chronic inflammation, leading to damage of the target tissue. Tissue debris is subsequently taken up by APC which have up-regulated expression of MHC class II and co-stimulatory molecules in response to inflammatory cytokines. These APC are then capable of activating CD4+ T cells specific for secondary tissue epitopes presented by the APC. The newly activated T cells then aid in destruction of the target tissue.\nDue to the requirement for presentation by APC of exogenous antigen, epitope spreading has historically been thought of as a phenomenon unique to CD4+ T cell responses. However, recent data have indicated that cross-priming by APC can participate in the induction of CD8+ cytotoxic T lymphocyte (CTL) responses as well. In particular, bone marrow chimera studies in murine tumour models have shown that tumour-specific CTL are predominately restricted to the MHC of the host rather than that of the tumour, suggesting that indirect presentation by host APC is involved in the generation of tumour-specific CTL. Moreover, there is increasing evidence that a pathway exists whereby exogenous antigen can be presented for eventual peptide loading onto class I MHC molecules. This phenomenon is best described for dendritic cells (DC) and provides a cellular mechanism to explain the process of cross-priming. Collectively, these data suggest that it may be possible for epitope spreading to occur during a class I MHC-restricted CTL response. Because re-presentation of MHC class I-restricted tumour antigens is known to occur, it has been postulated that if tumour-bearing hosts could initiate a CTL response against a single tumour antigen, that following tumor cell damage caused by the CTL, epitope spreading might occur via a mechanism analogous to that described in CD4+ T cell-mediated autoimmune diseases. Unlike during an autoimmune response, however, CTL epitope spreading during an anti-tumour response could be beneficial to the host by possibly allowing for elimination of variant tumor cells that have lost expression of the antigen (antigen negative tumour cells). [Markiewicz et al., (2001) International Immunology 13:625-632].\nMarkiewicz et al found that immunization with the single tumour peptide P1A followed by tumour rejection led to CTL activity against a P1A\u2212 tumour, indicating that the phenomenon of epitope spreading is not limited to CD4+ T cell responses. The population of CTL included cells recognized the unrelated antigen PIE. Since this epitope was not included in the vaccine and is a mutated peptide not presented in normal tissues, the source of PIE antigen must have been the tumor cell challenge.\nBecause many patients have ovarian tumours that express neither one of the aforementioned antigens there is a need to uncover additional antigenic targets for immunotherapy to manage localized and metastatic disease. Accordingly, provided herein are molecular targets for immunotherapeutic intervention in ovarian cancers.\nThe zona pellucida (ZP) forms an extracellular glycoprotein matrix surrounding the developing and ovulated oocyte and the preimplantation embryo and is also detectable in atretic follicles. The ZP induces acrosome reaction on sperm, determines the species specificity for fertilization and prevents polyspermy in mammals. The zona pellucida contains four major glycoproteins, ZP1, ZP2, ZP3 and ZP4. In vitro studies in mice indicate that O-linked oligosaccharide side chains of ZP3 are involved in the primary binding of the sperm to the ZP3, while ZP2 contributes to the subsequent and persistent ZP binding and functions as a secondary sperm receptor.\nThe ZP glycoproteins have been studied extensively for the development of vaccines for the fertility control of animals and humans. The proposed vaccine action is the induction in female subjects of effective sustained, but reversible levels of ZP-specific antibodies that inhibit sperm-egg binding and/or prevent sperm penetration of the ZP. Passive immunization of female mice with rat monoclonal antibodies against mouse ZP2 or ZP3 resulted in localization of the antibodies to intra-ovarian oocytes and long-lasting but reversible contraception. Active immunization of female mice with ZP3-derived peptides ZP3328-342, comprising a B-cell epitope recognized by the ZP3-specific contraceptive antibody, also led to reversible albeit incomplete contraception. These ZP3 peptides also induced a T cell response to the ZP3 peptide. These CD4+ ZP3 specific T cells adoptively transfer autoimmune ovarian disease (AOD) to syngeneic recipients. Since the desired contraceptive effect of ZP3 immunization is known to be mediated by antibodies, an acceptable contraceptive ZP vaccine should induce an adequate antibody response without activation of ZP3-specific T cells. Indeed a chimeric peptide consisting of a foreign T-cell epitope from bovine ribonuclease (RNase) and a minimal and modified murine ZP3335-342 B cell epitope has been designed that elicits antibodies to ZP and has a significant contraceptive effect without causing significant oophoritis/AOD. The bovine RNase T-cell epitope stimulates helper T cell (helper T lymphocytes, HTL) responses in mice, thus potentiating the contraceptive effectiveness without inducing ZP(3)-specific T cell action and T-cell mediated ovarian damage.\nImmunisation with (self)ZP antigen has also been used to study autoimmune ovarian disease (AOD). More in particular, animal models suitable for studying AOD have been reported wherein autoimmune disease was induced using ZP antigen vaccination. For example, it was demonstrated by Rhim et al. [Rhim et al. (1992) J. Clin. Invest. 89:28-35] that in B6AF1 mice T-cell and antibody response were induced by vaccination with mouse ZP3328-342 peptide. Further studies on truncated ZP3328-342 peptides substantiated that a T cell response is sufficient for induction of oophoritis; seven of such peptides lacking antibody binding sites, elicited severe oophoritis without concomitant antibody response. These peptides include a minimal oophoritogenic peptide of eight amino acids, ZP3330-337, which overlaps the seven amino acid antibody binding site, ZP3336-342, by two residues.\nIt was reported by Bagavant et al. [Bagavant et al. (1999) Biology of Reproduction 61:635-642] that transfer of ZP3 peptide-specific T-cells into na\u00efve recipient mice resulted in granulomatous oophoritis and enhanced ovarian expression of IL-1, TNF-\u03b1 and IFN-\u03b3. However the ovarian function of cell recipients was normal and the mice remained fertile. Antibody to ZP3 alone does not cause any ovarian pathology. Co-transfer of pathogenic T cells and ZP antibody together targets the inflammation into developing follicles leading to their destruction and the development of ovarian atrophy. In another study Bagavant et al [Bagavant et al. (2002) American Journal of Pathology 160:141-149] demonstrated that ZP3 peptide (human ZP3328-341, macaque ZP3328-341 and mouse ZP3330-342) immunization in primates can elicit a T-cell response and cause ovarian immunopathology that is similar to murine AOD.\nInternational patent application no. WO 2005/026735 (Buschmann et al.) relates to differentially expressed tumour-specific immunogenic membrane proteins and to their uses, in particular for finding at least one therapeutic molecule or compound which specifically regulates the expression of at least one of said membrane proteins, or for finding a therapeutic molecule that specifically binds to and/or interacts with any of said membrane proteins. The membrane protein can be SYPL, STOML2, RAGA, CLNS1A, PRNP, GNB2L1, GNG4, ITM2B, ITM1, TM9SF2, TM4SF6, OPRL1, LRP4, GLEPP1, TLR3 and/or ZP3. WO 2005/026735 teaches to administer the aforementioned therapeutic molecule or compound to neoplastic target cells for modulating proliferation, differentiation and/or cell migration of said neoplastic target cells. It is stated that the non-steroid dependent cancer to be treated results from the aberrant expression and/or biological activity of at least one of said immunogenic membrane proteins. It is also briefly mentioned in WO 2005/026735 that the development of a specific lesion, such as a pro-neoplastic lesion that can be found in epithelial tissues, into a neoplastic lesion can be inhibited by inoculating a subject with one of said membrane proteins adequate to produce antibody and/or T cell immune response. It is further specified that according to another embodiment the method comprises delivering one of the immunogenic membrane proteins via a vector directing expression of the said protein in vivo in order to induce such an immunological response to produce antibody to protect the subject from disease. Thus, WO 2005/026735, teaches several methods of suppressing the expression and/or biological activity of, amongst others, ZP3 membrane proteins in neoplastic target cells in order to modulate proliferation, differentiation and/or migration of said target cells, using either siRNA, receptor antagonists or antibody. WO 2005/026735 only discloses the expression of ZP3 membrane protein in certain colon cancer cells. No other reports of tumour associated expression of any ZP glycoprotein are known."}
-{"text": "A variety of therapeutic articles are used to treat various ailments, injuries and the like. However, a need exists for personal care articles, which can be used in home spa applications. The personal care articles need to be modular and have activatable personal care compositions to provide a variety of therapeutic benefits to a user."}
-{"text": "1. Field of Invention.\nThis invention relates generally to a medical ventilator and more particularly to a medical ventilator which is operable in both an invasive and non-invasive ventilatory environment.\n2. Description of the Related Art.\nVentilators are used by patients in various health situations. Typically, these patients have weak physiological attributes that prevent them from breathing unassisted and require invasive ventilation. Invasive ventilatory support generally requires the patient having either a tracheotomy or endotracheal tube disposed into the trachea of the patient. Such treatment generally occurs in hospitals and is administered in acute care situations or post operative recovery situations. Ventilators, such as the Siemens 300, are known to provide invasive ventilatory assistance.\nThe use of invasive ventilatory support has many risks. By positioning an endotracheal tube down the trachea of a patient, the patient is placed at some medical risk. Typically, there are certain physiological attributes in the trachea of a person for preventing bacteria and the like from invading into a person\"\"s lung. However, the endotracheal tube circumvents these natural defense systems. Accordingly, patients may acquire pneumonia by having the bacteria circumvent the natural defense mechanism of the body. Additionally, positioning of an endotracheal tube into a patient subjects the patient to a risk of tracheal abrasion. Overall, the positioning of an endotracheal tube, known as intubation, should be avoided when necessary.\nTypically, when a patient is placed on ventilatory support, most patients are subsequently weaned from the ventilator. Weaning involves manipulating the ventilator from that mode when the ventilator provides the most ventilatory assistance to the mode when the patient is providing most of the breathing. However, during the weaning of the patient from the ventilator, the patient remains intubated continuing further exposing the patient to health risks. Thus, there is a need for a medical ventilator that will enable a patient to be extubated, having the endotracheal tube removed from the trachea, as soon as possible to eliminate health risks.\nAdditionally, there are situations which arise wherein a patient having difficulty breathing is intubated prematurely and connected to a ventilator. This generally occurs since the attending physician lacks a ventilatory device which can provide ventilation in a non-invasive environment for initially determining if the patient is merely having difficulty breathing or truly requires invasive ventilation. This situation arises due to the costs associated with having a separate invasive ventilator and a non-invasive respirator which may be utilized to provide the patient with initial ventilation support. Due to the costs and expenses of having duplicate machines, most hospitals merely have invasive ventilators at their disposal.\nAlso, to provide ventilatory assistance in hospitals, the hospitals generally have specially designed respiratory care facilities having compressed air and oxygen hookups permanently affixed in a specific location. Such fully equipped facilities are expensive and also limit the areas where the ventilator may be accessed. Consequently, some patients who are otherwise healthy but require invasive ventilation are prevented from discharge due to their dependency on the respiratory care facility. With the high cost of hospital stays, some patients who require long term ventilatory utilization may occur exceptional hospitalization charges for the use of such expensive facilities. Thus, it is desirous to enable the ventilatory patient to be discharged to a low-acuity subacute facility or nursing home if the situation is warranted. However, most of these facilities lack the necessary pressurized air and oxygen hookups thus preventing the ventilatory patient from being discharged.\nAlso, since surgical areas and emergency room areas are typically stressful environments, it is desired that ventilators are easy to operate. This is also essential in today\"\"s health care environment since many different type of health care providers are assisting patients. These include respiratory therapists, nurses and physicians. Accordingly, it is desired that an intuitive ventilator exists for both invasive and non-invasive ventilation.\nAccordingly, it is an object of the present invention to provide a ventilator that can be used in both an invasive and non-invasive environment.\nFurthermore, it is an object of the present invention to provide a ventilator which can operate in both an invasive and non-invasive environment and having an operator interface that is simple to use for reducing error in operation;\nAlso, it is an object of the present invention to provide a self contained invasive/non-invasive ventilator having its own source of air and having the ability to mix with an oxygen or other gas source to provide flexibility in providing ventilatory assistance at different physical locations without requiring special respiratory care facilities with preexisting air sources.\nThe above objectives are accomplished according to the present invention by providing a ventilatory system for use in an invasive and non-invasive ventilator environment. The ventilator system includes a gas flow generator for providing a flow of gas to a patient. A conduit delivers the gas flow to the airway of the patient. At least a first valve regulates the delivery of the gas from the gas flow generator to the conduit. A controller controls the delivery of the gas flow to the patient. The system further includes a first set of operational parameters for directing the controller to control the delivery of gas to a patient if the patient is being ventilated in an invasive ventilation mode and a second set of operational parameters for directing the controller to control the delivery of gas to the patient if the patient is being ventilated in a non-invasive ventilation mode. Also a selector is utilized for selecting either the first or second set of parameters to direct the ventilator to provide either invasive or non-invasive ventilatory support to the patient.\nAlso, a unique blower is utilized for providing ventilatory assistance. The blower is a multi-stage centrifugal blower having an air inlet for receiving air from the ambient environment. A first impeller imparts centrifugal force onto the air. A first stator receives the air from the impeller and pressurizes the air. A second impeller subsequently receives the air from the first stator and imparts additional centrifugal force onto the air. A first impeller spacer directs the air from the first stator to the second impeller. A second stator receives the air from the second impeller and further pressurizes the air. A third impeller receives the air from the second stator and further imparts centrifugal force into the air. A second impeller spacer directs the air from the second stator to the third impeller. A blower outlet permits the pressurized air to leave the blower assembly. The overall impeller and stator configuration enables air to be pressurized to at least one hundred and sixty centimeters H2O when exiting the blower outlet.\nAlso, a graphical user interface is utilized in the ventilator for controlling the operation of the ventilator. An activation area is present for displaying a first activation device for a first mode of ventilation and a second activation device for a second mode of ventilation. A selector selects either the first or second activation device. A display area then displays the operational parameters pertaining either to the first or second mode of ventilation depending on the mode selected by the operator. Only those operational parameters relating to the particular ventilation mode selected by the operator are displayed.\nThese and other objects, features, and characteristics of the present invention, and the combination of parts and economies of manufacture, will become more apparent upon consideration of the following detailed description and appended claims with reference to the accompanying drawings, all of which form a part of the specification, wherein like reference numerals designate corresponding parts in the various FIGS."}
-{"text": "The present invention is directed to electrical amplifier apparatuses, and especially to electrical amplifier apparatuses adaptively controlling output to accommodate signal clipping to improve recovery from signal clipping.\nThe present invention is generally useful in electrical signal amplifiers. The present invention is particularly advantageously employed in audio signal amplifying apparatuses. Amplification of electrical signals, and especially audio signals, is a challenging area of analog circuit design. Unlike some other linear circuits, audio systems are required to function with high fidelity over a broad dynamic range. While the average power required for most music listening (an example of use of an audio amplifying apparatus) is quite low, there can be peak output levels which are larger than the supply rail\u2014the signal provided for operating the amplifier. When the output signal levels exceed the supply rail, clipping occurs at the outputs of the amplifier. Ideally the amplifier output signal will exactly match the shape of the input signal with a perfect cutoff for levels above the supply voltage. Class-AB power amplifiers typically will have no trouble achieving this kind of performance provided that good design techniques are employed. Class-D amplifiers however, which are rapidly increasing in popularity for audio amplifying applications, present challenging design issues in order to approach such ideal clipping performance.\nA problem that may be experienced during a clipping condition in an amplifier is that, while clipping (i.e., while the output signal levels of the amplifier exceed the supply rail), may be illustrated by way of example and not by way of limitation, in the contest of an exemplary Class D amplifier. The exemplary amplifier receives input signals at an integrator unit. The integrator unit integrates the input signals (using integrating capacitors in a feedback arrangement, which will be described in greater detail hereinafter in connection with FIGS. 1\u20134. When a clipping condition exists, the amplifier output can no longer maintain proportionality with the input signals received by the amplifier. During this clipping condition, integrating capacitors continue to be charged, so that the output of an integrating unit becomes larger than the reference signal. When the input signal returns to a lower level at which clipping does not occur, the amplifier takes a finite time to reduce the excess charge that has been accumulated in the integrating capacitors during clipping conditions. After the excess charge is reduced, the amplifier can return to normal operations. The delay in reducing excess charge in the integrating capacitors causes distortion during clipping recovery. A user of an audio amplifier system can notice noise occasioned by the distortion that occurs during clipping recovery operation.\nThere is a need for an amplifier apparatus that can reduce time required for clipping recovery operations."}
-{"text": "Many broadcast systems employ one or more headend sites (also known as headends) to receive broadcast signals from multiple sources and/or on multiple channels to facilitate redistribution of the broadcast signals to subscribers and/or other reception sites over a common transmission medium. For example, a typical Internet Protocol Television (IPTV) system employs at least one headend to receive Advanced Television Systems Committee (ATSC) broadcast signals on one or more broadcast channels for redistribution via a digital Internet Protocol (IP) network. The digital format of ATSC broadcasts, as well as the redistribution of their content via a digital IP network, offers the promise of improved signal quality over traditional National Television System Committee (NTSC) analog broadcasts. However, the quality of ATSC broadcast content, like other types of content conveyed by wireless broadcasting systems, can suffer degradation associated with various transmission phenomena or channel conditions, such as co-channel interference, adjacent channel interference and/or weather-induced signal fading."}
-{"text": "1. Field of the Invention\nThe present invention concerns a device to reduce the electromagnetic field load caused by the presence of a medical intervention apparatus in magnetic resonance examinations.\n2. Description of the Prior Art\nMany medical procedures today are conducted in an interventional manner. Diagnostic or therapy methods that\u2014in contrast to conservative procedures\u2014conduct targeted procedures (interventions) on pathological tissue in order to positively affect the course of the illness are designated as interventional. Invasive or, respectively, minimally invasive procedures are also examples of such methods, wherein medical instruments or implants are brought to the treatment or examination location by various techniques.\nParticularly in the treatment of vascular diseases, guide wires are used that are inserted by a physician into hollow vessels (lumens) of a patient. Later a catheter (for example) is inserted along this wire for diagnosis or therapy. The guide wire itself often has a diameter of less than 0.5 mm. Due to the high mechanical demands, such guide wires have been conventionally produced from metal and metal alloys. Magnetic resonance is a modality particularly suited for imaging of the human body with a high tissue contrast. It would be particularly suitable for observation of intervention apparatuses during their insertion and placement in the human body, but the electromagnetic fields in the two-digit or three-digit megahertz range that are required for imaging are problematical. Therefore, heretofore there have been high risks in the real-time tracking of guide wires and other electrically conductive intervention apparatuses given movement in the human body. This is highly risky because high-frequency currents can form on elongated electrically conductive apparatuses due to injection of the transmission fields. This leads to heating at corresponding powers and, in the disadvantageous case, to tissue burns inside the human body. The currents are even higher at such elongated intervention apparatuses if waveguide resonances occur along the entire length due to the intervention apparatus. If the lengths of the guide wires are also on the order of the wavelength of the magnetic resonance frequencies being used, the risk of endangering a patient is particularly high.\nA prediction about the potential for danger by means of simulation calculations is very difficult since the actual danger is situation-dependent. The length and orientation of the intervention apparatus in the patient, the conductive paths outside the patient, and even properties of the patient, play a role. Finally, the position of the intervention apparatus in relation to the transmission antenna, and the transmission antenna type, are also very influential.\nDue to the high risks from the induced radio-frequency currents, the real-time tracking of intervention apparatuses with the aid of magnetic resonance has previously been realized with the transmission power being drastically limited. Another solution is the construction of guide wires or catheters from non-conductive materials. However, all desired material properties cannot be ensured with these materials; such intervention apparatuses are additionally very expensive and complicated."}
-{"text": "The present invention relates to amplifiers and more particular to amplifiers employing shunt feedback.\nThe traditional series shunt feedback architecture illustrated in FIG. 1 has been widely hailed as the preferred architecture for Low Noise Amplifiers (LNAs) requiring a well defined input impedance. However combining series shunt feedback with Gilbert\"\"s current steering techniques for providing an Automatic Gain Control (AGC) function presents certain difficulties.\nThe following equations describe how the input impedance Rin of the traditional series shunt feedback amplifier is set by a negative feedback amplifier:\nRf=Rs*(1+A)\nRin=Rf/(1+A)\nwhere Rf is the effective feedback impedance, Rs is the source impedance, and |A| is the magnitude of the gain of the amplifier. It will be appreciated that as long as Rf remains constant, then so does the input impedance Rin. However, if the amplifier has a varying gain, then the input impedance Rin will also change.\nA potential solution to this problem is offered by the architecture of FIG. 2, where the input impedance of an amplifier arrangement is set by an input stage comprising an amplifier which has a fixed gain. The output of the input stage feeds a variable gain amplifier stage. This solution is not optimal however, since the input stage has a fixed gain and the input signal to the variable gain amplifier stage is therefore higher than the input signal to the amplifier arrangement as a whole. This creates a problem in signal handling requirements and may result in a power in-efficient solution.\nIn view of the disadvantages noted above, existing amplifier arrangements may not meet the strict requirements of modem electronic systems, such as interactive cable television equipment, in which a range of high performance RF products are required ranging from mixer oscillator circuits to LNAs with AGC and loop through circuits.\nIt is an object of the present invention to provide an amplifier arrangement having a stable input impedance. It is a further object of the present invention to provide a low noise amplifier which is power efficient. These and other objects are achieved at least in part by an amplifier arrangement in which the amplifier stages thereof are operated in current mode.\nAccording to a first aspect of the present invention there is provided an AC voltage amplifier arrangement comprising:\nan input stage having an input, an amplifier connected in series shunt feedback mode, and an output; and\na second stage having an input coupled to the output of said input stage, and an output,\nwherein said input and second stages are arranged to operate in current mode.\nPreferably, said second stage is a variable gain amplification stage.\nPreferably, the voltage amplifier arrangement is arranged to amplify radio frequency voltage signals.\nPreferably, said amplifier of the input stage is a transconductance amplifier. More preferably, the transconductance amplifier comprises a bipolar junction transistor configured as a common emitter stage, or a differential pair of bipolar junction transistors configured as common emitter stages.\nPreferably, the arrangement comprises a unity gain buffer amplifier coupled in the series shunt feedback loop of the input stage.\nPreferably, said second stage comprises a current steering arrangement arranged in use to steer a proportion of the current, provided to the input of the second stage, to an output load, whereby an output voltage is developed across said load and is available at said output of the second stage. The current steering arrangement may be a Gilbert steering stage. More preferably, said current steering arrangement comprises a pair of bipolar junction transistors having their emitters coupled together and to said output of the input stage. The current steering arrangement may comprise two pairs of bipolar junction transistors in a differential arrangement, coupled respectively to differential outputs of the input stage.\nThe second stage may be arranged to be controlled by an Automatic Gain Control (AGC) signal in order to maintain the output of the stage at a desired level. More preferably, where the stage comprises a pair (or pairs) of bipolar junction transistors, said AGC signal is applied to the bases of the transistors.\nThe second stage may comprise a mixer or a buffer.\nAccording to a second aspect of the present invention there is provided a method of amplifying an AC voltage signal, the method comprising:\nconverting the voltage signal into a current signal;\namplifying said current signal to a desired level; and\nconverting the amplified current signal into an amplified voltage signal.\nPreferably, the amplification factor by which the current signal is amplified is less than one.\nPreferably, amplification of the current signal comprises splitting the current signal into at least two parts, with one of said parts being converted into the amplified voltage signal."}
-{"text": "Social media platforms are becoming an increasingly important way to interact and engage with people, such as customers of a business. Both customers and businesses can benefit from an increase in the number of opportunities to engage through the use of social media. In particular, the benefits of engaging are particularly pronounced when a customer attends an event at the physical, geographic location of a business or event. The ability to recognize the customer's physical presence and engage with the customer in real time, can be used to enhance the customer's experience at the location, and thus is extremely valuable.\nHowever, some techniques for determining the presence of a customer at a physical location using social media are extremely limited in their ability to effectively determine relevant location information, particularly in real time. For example, some existing techniques rely solely on social media posts with geotagged physical locations in order to determine social media activity by customers that are physically present at a location. However, techniques that rely solely on geotagged social media posts can miss the vast majority (e.g., up to about 95%) of social media activity relating to the geographic location of interest."}
-{"text": "1. Field of the Invention\nThe present invention relates to computer networks and in particular to handling client requests in systems having multiple servers.\n2. Description of Related Art\nOn computer networks, nodes may interact through requests for data, services, or other resources. Client nodes may generate requests and server nodes may service those requests. Nodes may be stand-alone computers, servers, or other computing devices, as well as virtual machines, threads, processes, or other network participants. To handle a greater volume of requests, a plurality of servers may be allocated to responding to one or more types of requests and the servers as a group may be referred to as a server farm. For example, a bank, ticket office, or some other organization may use a server farm to process requests from clients. The servers may be web servers on the Internet responding to web clients.\nTo provide coherent service to clients, state or session information may be collected and used by servers, so that their response to a given client may depend on session data for the client in the server farm. For example, session information may be used to implement a \u201cshopping cart\u201d used across multiple requests over one or more servers. It may be more efficient to store the session information on a single server within the server farm, rather than disperse that information, and ensure that the corresponding server services the requests from the associated client.\nA typical system may have a unique session ID assigned to correspond to collected client session information. Initially, requests do not contain session IDs since there is not yet a session formed with the server farm. The first server to receive a request from a given client typically generates the session I). After a server creates a session for a client, the server generates a session ID to uniquely represent that session, and instructs the client to return the generated session ID in subsequent requests from that client. For example, a session ID may be encoded in the URLs the client is instructed to access or it may be stored in a server-generated cookie the client is instructed to send back to the server with each request.\nRequests may be routed between the clients and the servers through one or more devices, which may be considered upstream from the servers and downstream from the clients. These devices may be part of the server farm or may operate within the network or other system. For example, the server farm may include one or more routers that receive the requests from the network and route them to servers.\nThe work that a server performs processing requests it receives is referred to as the load of that server. Depending on the allocation of requests and the processing required, the load may differ widely between servers and limit performance. Another example of upstream devices is load balancers. In order to distribute load across a plurality of servers within a server farm, load balancers may be used. Load balancers may be additional devices, may be part of other devices, or may include other functionality. For example, load balancers may receive incoming requests to the server farm and route the requests to individual servers in addition to determining where requests may be sent to better balance loads in the server farm.\nIn some systems, a load balancer or other device may be configured to receive incoming requests, examine them for a session ID, and then route them to the server that generated the session ID so that session information is kept on a single server within the server farm. The technique of using the session ID to select the server that will service a given request is called \u201csticky sessions\u201d, \u201csticky load balancing\u201d, or \u201csticky routing\u201d since the presence of a session ID causes the requests associated with that session to \u201cstick\u201d to a specific server.\nSome systems may include load balancers or other devices in the server farm that receive requests and send them to servers. These devices may maintain a mapping of session IDs to servers. These devices may then compare the session IDs with the requests to their mapping information to find the corresponding server. However, the size of the data for the mapping may be prohibitive as there may be many session IDs-at least one per client of the server farm.\nIn some systems, servers may provide their address, such as an Internet hostname or address, to clients so that the client or devices within the server farm, such as load balancers, may use that information to route subsequent requests using standard Internet routing protocols to the corresponding server. However, exposing the server addresses or other details of the sever farm configuration communicates information about the internal arrangement of the server farm to the client or others outside the server farm. As the client is typically a distrusted third party, exposing information about the internal operation of the server farm may be a security risk.\nAlternatively, an administrator may configure a load balancer and servers with a particular ID corresponding to each server as identifying information. Servers may then include this identifying information in their session ID without revealing information about the internal structure of the server farm. However, the load balancer and server must then share data that allows the load balancer to target a specific server from the identifying information. This implementation imposes a burden on server farm administrators, as configuration information for these session IDs is then maintained for all load balancers and servers."}
-{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor memory device and a manufacturing method of the semiconductor memory device.\n2. Related Art\nData reading from a NAND flash memory is performed in a page unit. One page includes 4 k-bit to 64 k-bit data, for example. Therefore, in read operation, data of 4,000 to 64,000 memory cells are read simultaneously. Because cell current flows simultaneously to the 4,000 to 64,000 memory cells at this time, potentials of cell source lines sometimes increase (or is lifted). When the potentials of the cell source lines increase, cell current decreases. Consequently, a threshold voltage of the memory cells increases, and a distribution of the threshold voltage of the memory cells is spread. To suppress the spreading of the threshold voltage distribution, the cell source lines connecting between source drivers and sources of the memory cells preferably have as low resistance as possible. For this purpose, a top-layer wiring layer of low resistance in a multilayer wiring structure on a memory cell array can be used as cell source lines and cell well lines (see Japanese Patent Application Laid-open No. 2005-142493 and Japanese Patent Application Laid-open No. 2006-245547).\nIn recent years, to improve performance of memories, the ABL (All Bit Line) sense scheme has come to be used frequently. In the ABL sense scheme, sense amplifiers are provided corresponding to all bit lines and all bit lines are simultaneously read. Therefore, a large amount of current is consumed in the sense amplifiers. To supply power to the sense amplifiers, wide power source lines need to be arranged from a pad to the sense amplifiers. This increases chip size.\nTo further downscale the chip size, there is considered a layout of forming on the top-layer wiring layer of the memory cell array a VSS (ground potential) power source lines to be connected to the sense amplifiers, as well as the cell source lines and the cell well lines (see Kanda et al. \u201cA 120 mm2 16 Gb 4-MLC NAND Flash Memory with 43 nm CMOS Technology\u201d ISSCC 2008/SESSION 23/NON-VOLATILE MEMORY/23.6 p.p. 430, 431 and 625). To decrease resistance of the wiring formed on the top-layer wiring layer in this way, an installation area of each wiring is preferably increased. To increase each wiring area of the cell source lines, the cell well lines, and power source lines as far as possible, each wiring formed on the top-layer wiring layer can cover a whole upper part of the memory cell array. However, when the wiring covers the whole upper part of the memory cell array, hydrogen is not thoroughly supplied to the memory cells in a forming gas/anneal process performed at a final stage of a memory-chip manufacturing process. This becomes a cause of degradation of a data holding characteristic of the memory cells."}
-{"text": "The subject invention relates to interelectrode separators for eliminating zinc dendrite bridging between opposite polarity electrodes in secondary alkaline zinc batteries.\nBattery separators function essentially to electrically isolate the positive and negative plates from each other. They prevent direct contact and suppress \"treeing\" or interelectrode dendrite growth which causes shorting of the respective plates. An ideal separator would isolate the plates without inhibiting electrolyte mobility, and without increasing the battery's internal resistance.\nFor many years, regenerated cellulose was the most commonly used separator for secondary zinc alkaline batteries (e.g., silver-zinc, nickel-zinc). In this regard, fibrous sausage casing and unplasticized cellophane (i.e., 0.025-0.5 mm dry thickness) were wrapped at least one time around the zinc electrode. Often several turns of the separator were used to insure adequate dendrite resistance. These cellulosic separator materials were generally utilized in conjunction with an \"interseparator\" which was a loosely woven or bibulous material wrapped or heat-sealed around the positive electrode to prevent the cellulosic separator from contacting the highly oxidizing environment of the positive electrode. Cellulosics performed well in vented, flooded cells but lacked long-term stability in sealed cells and at temperatures above the ambient (i.e., up to 100.degree. C.). Other problems with cellulosic separators involved swelling and degredation due to oxidation.\nMore recently, newer types of separators with better chemical stability and physical properties have become available. One such class of separators has been designated I/O inorganic/organic) separators and comprises a major inorganic component in a highly stable organic polymer binder. Such separators typically include various inorganic oxides, hydroxides, etc. in a fibrillated polytetrafluoroethylene (hereafter PTFE) binder. These materials are highly flexible and quite durable in the alkaline electrolyte of a secondary zinc battery. The process for making such materials essentially consists of the mixing of the PTFE and inorganic materials under mild shear conditions (e.g., liquid blender) with a high-boiling naphtha lubricant at temperatures above 19.degree. C. to cause fibrillation of the PTFE (see HS-10 bulletin by E. I. Du Pont de Nemours & Co., Plastics Dept., 1968). PTFE levels are normally in the 1-25% by weight range.\nIt is an object of the present invention to provide a new, low cost separator which is strong, flexible and particularly useful in preventing zinc dendrite growth between the electrodes of a secondary alkaline zinc battery. This and other objects and advantages of the present invention will become more readily apparent and understood from the following detailed description and illustration thereof."}
-{"text": "This invention relates to a method of depositing a high-emissivity carbon coating on a substrate by reactive RF sputtering. More specifically, the invention relates to a method of depositing a carbide-carbon coating on a substrate by reactive RF sputtering having a composition gradient such that the coating is composed primarily of a carbide compound near the substrate and primarily of carbon near the exterior surface.\nThermoelectric generators are frequently used as an energy source in aerospace applications. Temperature regulation in these generators is a critical problem for both orbital satellites and deep space probes. A high-emissivity coating on the metallic thermoelectric generator casings would improve the performance of satellite thermoelectric systems by radiating large quantities of heat while operating at low temperatures. The high-emissivity coating must be capable of withstanding the stresses imposed during launch and the elevated temperatures and vacuum conditions encountered during operation.\nCarbon has many of these properties and would serve well as a high-emissivity coating. However, carbon does not adhere well to the metals often used in thermoelectric generator casings, namely iridium and aluminum. Attempts have been made to blend carbon with an epoxy and paint the mixture onto the casings, but this process adds undesirable weight to the generator. A more promising method is to sputter deposit a coating with a composition gradient such that the surface of the coating near the substrate is rich in material that adheres well to the substrate while the exterior surface of the coating is rich in carbon. Previous methods of deposition have involved rotating a substrate between two targets of different materials and gradually changing the deposition parameters on the two targets to obtain the desired composition gradient. This method is too complicated to be practical for the massive parts of the generator which must be coated."}
-{"text": "1. Field of the Invention\nThis invention generally relates to the measurement of a multiple phase flow. Here, microwave radiation is used to measure the two-phase fluid content of oil field flowlines and also the two-phase fluid content of fluid-saturated porous media.\n2. Background Information\nA multiple fluid mixture has at least one physically distinct and mechanically separable fluid within at least one other physically distinct and mechanically separable fluid in a physical/chemical system. In the case at hand a multiple fluid mixture would be something such as a combination of water and oil.\nSome of the factors that complicate the measurement of the multiple fluid mixture are: whether the mixture is in a static or dynamic state; and whether the method is intrusive or non-intrusive. One simple method of analyzing a multiple fluid mixture involves capturing of a sample volume of the mixture and then analyzing it (or simply allowing the mixture to separate into its component phases). This method, however, is an intrusive one and involves too much time and effort to capture a sample volume. Furthermore, when captured, there is no guarantee that the sample volume is an adequate representation of the entire sample. If so, the calculation for the total concentration of the two components in the mixture will be erroneous. For that reason a non-intrusive method would be preferable.\nMultiple fluid mixtures may also be measured by subjecting a sample (either in a static or dynamic state) to photonic laser, or infra-red radiation. In these cases, a multiple fluid mixture is subjected to the specific radiation and the attenuation of the radiation is used to calculate the component phases. See U.S. Pat. Nos. 4,157,470 to Kotuka et al.; 4,470,647 to Chraphyvy et al.; and 4,394,575 to Nelson; see also V. E. Shrock, \"Radiation Attenuation Techniques in Two-Phase Flow Measurements\", Nat'l ASME/AICHE, Heat Transfer Conference, Minneapolis (1969); and Davis, \"VHF Electrical Measurement of Saturations in Laboratory Corefloods\", SPE 8847 (1980).\nThe other factor that was mentioned above, i.e., whether the mixture is in a static or a dynamic state, would reflect mainly on the particulars of the measurement system. For example, in most applications it is most desirable to measure the mixture when it is in a fluid or flowing state. A common example would be an oil field flowline pipe that may transport liquid. However, there are examples where it may be desirable to measure multiple fluids in a static system. For example, a laboratory situation may be set up in which a sample petroleum formation, containing a porous medium (such as Berea Sandstone), has a two-phase mixture. Both cases could be applications of the same measurement technique.\nThe techniques described above are useful in understanding how fluid concentrations may be measured in a multiple fluid mixture. However, as applied to petroleum applications, i.e., multiple fluid mixtures of water and oil in a flowline or a test core, microwaves work to distinguish between multiple fluids. Some work has already been done in which microwaves are used to determine the concentration of multiple fluids, as outlined below. The devices disclosed in U.S. Pat. Nos. 4,167,736 to Tomlinson; 4,196,385 to Vesterpaard et al.; 4,423,623 to Ho et al.; and 4,503,384 to Nagy et al. all are intrusive into the fluid that they are measuring and are undesirable because of that fact. U.S. Pat. Nos. 4,289,020 and 4,301,400, to Paap, show devices that manipulate the fluid channel through which crude oil flows and are undesirable because they contain \"dead spaces\" where the flow may eddy and throw off the measurement system. U.S. Pat. No. 3,586,971 to Besido shows a device that monitors the moisture content in the sheets of material. It uses a method that detects a resonant frequency shift only. U.S. Pat. No. 3,818,333, to Walker, shows a microwave window and antenna, but does not disclose how measurements are made to determine fluid content. The following patents are issued to Lorne A. Davis et al. U.S. Pat. No. 4,519,982 discloses a device that scans a two-dimensional sample. U.S. Pat. Nos. 4,482,634 and 4,490,676 disclose a method ('634) and an apparatus ('676) for monitoring fluid in a core of material. However, both patents only measure power level, not phase shift. U.S. Pat. Nos. 4,486,714 and 4,543,821 use microwave circuitry that is similar to 4,482,634, but show different methods of fluid delivery to a test core. Two articles relate to some of these patents and are: Haskin and Davis, \"A Comparison of Laboratory Linear and Pattern Flow Chemical Floods Using a Volumetric Linear Scaling Concept for Oil Saturation and Distributions\", SPE 10197 (1981) (see U.S. Pat. No. 4,519,982); Brost and Davis, \"Determination of Oil Saturation and Distributions in Field Cores by a Microwave Spectroscopy\", SPE 10110 (1981) (see U.S. Pat. Nos. 4,482,634 and 4,490,676).\nOther articles disclose the use of microwaves to measure fluid contents. R. W. Parsons was one of the first investigators to disclose a method to determine in situ saturations of multiple fluid phases in a laboratory coreflood using microwave attenuation, see Parsons, \"Microwave Attenuation--A New Tool for Monitoring Situations in Laboratory Flooding Experiments\", Soc. Pet. Eng. Jour. 302 (August, 1975). However, the signal penetration depth in this method is limited to one inch and the frequency is 21.26 Gigahertz. P. E. Baker discusses a similar method (at the end of the Parsons' paper, see page 309). Bentsen and Saeedi show a method that is similar to Parsons, see Bentsen and Saeedi, \"Liquid-Immiscible Displacement in Unconsolidated Porous Media\", Jour. Can. Pet. Tech. 93 (January-March, 1981). Their method measures only microwave attenuation, has a penetration depth of only 1.27 centimeters, and operates on a microwave frequency of 27 Gigahertz. Wasan et al. also show a method similar to Parsons, see Wasan et al., \"Microwave Spectroscopy Analysis of Surfactant/Polymer Flooding: Interrelationships Between Chemical Slug Properties, Coalescence Phenomena, and Tertiary Oil Recovery\", SPE 8327 (1979). For example, Wasan et al. measure microwave attenuation using a frequency of 21.83 Gigahertz. Another similar device is described in Gladfelter and Gupta, \"The Effect of Fractional Flow Hysteresis on Recovery of Tertiary Oil\", Soc. Pet. Eng. Jour. 509 (December, 1980). Here, measurements of multiple fluids are made with microwave attenuation and depth penetration is limited to 1.91 centimeters. An article that reviews all the current methods to measure multiple fluids and petroleum applications is Saraf, \"Methods of In-situ Saturation Determination During Core Tests Involving Multiphase Fluid Flow\", Report 1981-6 Petroleum Recovery Institute (March, 1981).\nThese references mostly use only one measurement technique, i.e., signal attenuation, to determine fluid concentrations. This has some drawbacks, in that: some physical properties are masked in different flow regimes when you use signal attenuation; you may be limited to determining the concentrations of only two fluid phases (with one measurement technique); and you do not have a confirmation of the fluid concentrations when you use a single method."}
-{"text": "The invention relates to a battery terminal connecting cable with a strand consisting of numerous fine wires and a further section having a hole for a screw, whereby on the further section the numerous fine wires are welded together, a method for the manufacture of a battery terminal connecting cable and a device for the manufacture of such a cable.\nBattery terminal connecting cables are used primarily to interconnect battery cells conductively. In this case one also talks of cell connectors. For the conduction of relatively high currents at low electrical resistance the strand usually consists of copper wires twisted into a strand having a cross-section of approximately 50 mm2. Depending on the field of application, however, both thicker and thinner strands are used.\nIn order to connect such a strand consisting of numerous fine copper wires to a battery terminal, a copper pipe section is first inverted over the strand and this pipe section is then pressed to an approximate plate shape. In this plate there is a hole which first passes through the upper side of the original pipe section, then through the compressed cable and finally through the lower side of the original pipe section. Finally in this hole there is inserted a screw which interacts with a thread in the battery terminal so that when the screw is tightened, the strand held together by means of the copper sleeving is pressed onto the battery terminal.\nIn practice it has been found that even a relatively strongly tightened screw works loose as a result of vibrations, as occur especially in batteries located in vehicles. This has the result that the end of the cable is no longer securely connected to the battery terminal. Contact is therefore made over smaller areas and is thus severely impaired. In cases of more severe loosening this leads to undesirable heating, even to the extent of sparking.\nIn order to prevent screws working loose, in known battery terminals a plastic section, such as preferably a plastic bead, is provided in the thread region between the screw and the battery terminal. This bead becomes deformed as the screw is turned and acts to prevent loosening of the screw.\nThe insertion of a bead or other plastic section has the result that expensive, special screws are required. Since battery terminal connectors are a mass-produced product manufactured in large quantities, any increase in the cost of the product leads to economic disadvantages.\nA battery terminal connecting cable of appropriate type is known from EP-A-O 707 321. In practice the end of the cable strand produced by spot welding does not give the required strength values for the cable end so that the cables described have not proved suitable as battery terminal cables.\nU.S. Pat. No. 4,325,760 discloses a cable whose end was soldered. However, soldering is an expensive procedure which leads to increased costs in the area of mass production.\nThe invention is based on the problem of preparing a connection between battery terminal connecting cable and battery terminal at a favourable cost, which ensures secure contact between the battery terminal connecting cable and the battery terminal even under severe vibration or thermal expansion.\nThe problem is solved using a battery terminal connecting cable of appropriate type where the strand is at least partly insulated and a further section is secured at the end of the strand.\nThe invention is based on the knowledge that the problem of screw loosening in the battery terminal is not primarily attributable to the screw being held inadequately in the battery terminal. The reason for loosening screws lies in the fact that at times fairly high temperatures occur at the contact piece of the battery terminal connecting cable as a result of the flowing currents. At varying temperatures the fine wires inside the copper sleeving are severely compressed from time to time and then remain in this severely compressed form. As a result after many temperature fluctuations the originally tightened screw is subject to less force so that the screw sits more loosely in the thread. Loosening of the screw can then only be impeded by the plastic section described above. The plastic section certainly impedes loosening of the screw but the fact that the screw at times exerts less pressure on the contact piece is not impeded by the plastic section. However, a constant pressure on the contact piece is required to ensure a continuous equally good current flow.\nPractical tests have shown that when a cable end consisting of wires welded together, is screwed on, the compression of the connecting piece described above does not occur and consequently normal screws without protection against loosening can be used. Since the screw now presses the contact piece onto the battery terminal at a constant pressure, optimal electrical conditions are created and thus the heating in the contact region is reduced.\nAccording to the invention a further section is secured at the end the strand. This further section can, for example, serve as a contact piece and can be connected in various ways to the end of the strand welded to form a solid material.\nA particularly preferred embodiment of a battery terminal connecting cable is obtained if the further section is welded onto the end of the strand. In this way a specially favourably shaped contact piece can be used. This in particular allows the weight of the battery terminal connecting cable to be reduced.\nIt is also advantageous if the further section, preferably constructed as a contact piece, is made of copper.\nIn all known battery terminal connecting cables the longitudinal axis of the further section extends in the direction of the axis of the battery terminal connecting cable. Since the battery terminal connecting cable serves as a connecting element between two terminals, it seems to be logical to use a further section extending in the direction of the axis of the battery terminal connecting cable. However, tests have shown that in many cases it is of great advantage if this further section is configured such that its longitudinal extension runs at an angle to the longitudinal axis of the cable. Examples of this type of configuration are shown in FIGS. 6-9.\nDepending on the configuration of the battery terminals, in practice the connecting cables are frequently very severely bent. In many cases, the angular configuration of the further section to the axis of the battery terminal connecting cable reduces the necessary bending of the cable and also allows shorter cables to be used. The configuration described thus has the result that the cable can be protected and material usage can be reduced.\nThe angular configuration described is also of great importance for battery terminal connecting cables regardless of the afore-mentioned characteristics and can thus also be used to advantage for clamp connections between the cable strand and a further section.\nThe problem according to the invention is also solved using a method for the manufacture of a battery terminal connecting cable in which a strand consisting of numerous fine wires is welded by means of ultrasound and the strand is welded to a contact piece or with a contact piece.\nThe use of ultrasound for making electrical connections is already known from EP-A-0 671 790. However, the welding of the strand to a contact piece or with a contact piece cannot be inferred from this document.\nAll types of welding or soldering of an end of a strand developed so far have failed in that the high heat input during the welding process is transferred from the copper cables normally used to the insulating layer. On the one hand, this leads to destruction of the insulation and on the other, the input heat is rapidly dissipated. Only welding by means of ultrasound produced such good results that the insulating layer can remain on the cable strand during welding. This surprising result reduces the energy input and allows known manufacturing methods for moulding plastic onto the connecting regions to be used even for the new battery terminal connecting cable.\nFor manufacture of the battery terminal connecting cable according to the invention a device is proposed which has a sonotrode, a strand feed device, a contact piece feed device and a pressure cylinder arranged so that the strand and the contact piece can be pressed onto one another by means of this device.\nThe device allows the battery terminal connecting cable to be manufactured automatically and rapidly.\nSince during pressing of the wire strand the fine wires are pressed to the side, it is proposed that on the device according to the invention there should preferably be provided movable jaws which act on the strand at right angles to the axis of the pressure cylinder. These jaws together with the pressure cylinder and the sonotrode define a cross-section in which the wires should be welded together to form the most compact copper piece. At the end of the welding process the jaws are moved away from the strand and the pressure cylinder is retracted so that the strand connected to the contact piece can be removed from the device."}
-{"text": "1. Technical Field\nThe present invention relates to a method of upgrading low-value linear alpha olefins (LAOS) of C18 or more. More particularly, the present invention relates to a method of manufacturing high-quality lube base oil by converting a low-value LAO having a C18 or more hydrocarbon chain into a linear internal-olefin (LIO) having a controlled double bond position distribution, which then undergoes oligomerization.\n2. Description of the Related Art\nEthylene is a raw material used in the chemical industry to the extent that the production and consumption thereof are considered as indicators of the chemical industry scale of a country. Typically, ethylene is employed as a monomer for preparing a polymer such as polyethylene. In some cases, LAOs having a C4-C40 carbon length (or chain) are prepared by adjusting the degree of polymerization of ethylene, and thus utilized to manufacture a variety of chemicals. In this regard, as used herein, the term \u201cfull range LAO\u201d refers to a C4-C40 LAO mixture resulting from ethylene polymerization or oligomerization.\nThe reason why LAO production techniques are important is that LAO is a chemical without sulfur and nitrogen derived from crude oil. Generally, since crude oil contains impurities such as sulfur and nitrogen components in amounts of ones of wt %, it is difficult to directly prepare chemicals composed of pure hydrocarbons without such impurities. However, ethylene resulting from catalytic cracking of crude oil is converted into LAO, which is then converted again into a desired chemical, thereby obtaining chemicals composed of pure hydrocarbons without impurities.\nEthylene polymerization is mainly carried out using a batch reactor under air-sensitive conditions in the presence of a metallocene catalyst. Since a metallocene catalyst has a very strong single acid site, ethylene is selectively polymerized at the single acid site of the catalyst, and is thus linearly grown at the acid site of the catalyst. Ultimately, the olefin produced by the LAO preparation reaction possesses a linear alpha-olefin structure having an even number of carbon atoms because of the use of an ethylene monomer.\nWhen operating conditions for the LAO preparation reaction are controlled, the LAO composition distribution may be altered. For example, when the reaction temperature or the polymerization time is increased, the proportion of LAO having a relatively high molecular weight may become high. In contrast, when the reaction temperature or the polymerization time is decreased, the proportion of LAO having a relatively low molecular weight may increase. Taking into consideration the industry scale for LAO-derived chemicals by LAO manufacturers, the LAO composition is changed in such a manner that the reaction conditions are appropriately controlled, thereby ensuring optimal profitability for individual projects.\nAlpha-olefins produced by the LAO preparation processes exhibit properties varying depending on the number of carbon atoms thereof, and the properties of the resulting chemicals may become different depending on the type of LAO feed. For example, C4 LAOs obtained by polymerizing two ethylene monomers are present in a gas phase, and polymers resulting from the polymerization thereof have an excess of intramolecular branches, making it difficult to apply to a specific product, for example, lube base oil. Also, even when C6 LAOs obtained by polymerizing three ethylene monomers are converted into polymers, such polymers are unlikely to be applied to lube base oil and so on due to many intramolecular branches thereof. As such, copolymerization of ethylene with a C6 LAO may result in an ethylene-1-hexene copolymer having properties different from those of conventional polyethylene. Also, C8 LAOs obtained by polymerizing four ethylene monomers may be applied to Group IV lube base oil via polymerization, and may be used to produce a copolymer with ethylene, as in the C6 LAOs. Also, C10-C12 LAOs are polymerized and thus mainly used as Group IV base oil, and C14-C16 LAOs are reacted with amine or succinic acid, so that the resulting products may be applied to various functional chemicals, or mixed and thus used for cheap drilling oil.\nThe aforementioned LAOs have certain end-uses, and the markets adapted for the demand thereof are established. Furthermore, the related techniques have been already developed to some extent. The most expensive LAO is a C8-C12 LAO for preparing high-quality lube base oil (e.g. Group IV lube base oil), and C6 and C8 LAOs may be mainly utilized because the demand thereof is high for producing copolymers. Also, although C14-C16 LAOs are relatively less used due to limited end-uses thereof compared to the C6-C12 LAOS, C14-C16 LAOs may be used to prepare specialty chemicals and may thus be formed into products by increasing purity thereof. However, C18 or more LAOs are problematic because the end-uses thereof are limited, and the purity thereof has to be increased to obtain desired products. In this regard, C18 or more LAOs have been employed as an additive for lube oil or in wax form, but are not general-purpose chemical products, and thus the demand thereof is undesirably low. Also, C20 or more LAOs are limited in serving as feeds because they exist as a solid at room temperature. Like this, the market prices of C18 or more LAOs are maintained very low compared to the other LAOS, on account of low demands, tendencies to exist as a solid, and restrictions where such feeds may be prepared into products only when the purity thereof is increased.\nOwing to the limitations of end-uses as above, the yield of high-value C8-C12 LAOs is intended to be selectively increased while decreasing the yield of C18 or more LAOs, by LAO manufacturers. However, the LAO products have a composition of the normal distribution due to characteristics of the polymerization catalyst, making it impossible to attain operating conditions for preparing only C8-C12 LAOs and for drastically lowering the yield of C18 or more LAOS.\nHence, LAO preparation under reaction conditions able to lower the yield of C18 or more LAOs as much possible is chosen by LAO manufacturers, and is advantageous because only LAOS, the demand of which is high and the value of which is also high, may be selectively prepared while decreasing the yield of C18 or more LAOs regarded as byproducts. However, since the composition of products resulting from ethylene polymerization has the normal distribution as mentioned above, the yield of C8-C12 LAOS, corresponding to feeds for end-uses (especially Group IV lube base oil) which have the highest price and a future market of which is expected to grow, is relatively low, which is undesirable.\nAlternatively, the reaction conditions may be set so as to maximize the yield of the most expensive C8-C12 LAOs without taking into account the yield of C18 or more LAOS. As such, the yield of C18 or more LAOs is 10 wt % or more, but the yield of relatively high-value C8-C12 LAOs may be increased. Therefore, if C18 or more LAOs obtained by full range LAO techniques can be simply converted into general-purpose high-value products, the effect of such technology is expected to be very significant. Accordingly, methods of converting C18 or more LAOs into high-value chemicals are required.\nMeanwhile, techniques for preparing chemicals (e.g. drilling oil) having improved properties in such a manner that a double bond present at the alpha position of the LAO is shifted to the inside of the carbon chain to form a linear internal-olefin (LIO) (olefin shift reaction) are also known (e.g. U.S. Pat. No. 6,281,404). However, the above techniques are not regarded as the application of low-value C18 or more LAOs because a wide range of C4-C50 LAOs are utilized, and also because the double bond shift reaction for 1-pentene is mainly disclosed. In addition, there are known techniques for manufacturing lube base oil (e.g. EP 0 136 377 B1) by oligomerizing internal olefins (especially, internal olefin produced from dehydrogenation of paraffin so that a double bond is randomly distributed in the molecule) in the presence of a boron trifluoride catalyst and a promoter, and for manufacturing lube base oil (Japanese Patent Application Publication NO. 1993-306241) by oligomerizing a C10-C24 LAO/LIO mixture. However, these techniques are limited in improving properties of the lube base oil because the position of the double bond in the olefin is not controlled. In particular, such techniques are unlikely to become technology for increasing the value of low-value C18 or more LAOS produced inevitably as byproducts in full range LAO preparation."}
-{"text": "The present invention relates generally to an optical device, and more particularly, to an optical device utilizing a superlattice of one or more indirect bandgap materials.\nThe strengths of optical transitions between conduction band states and valence band states in semiconductors are greatly influenced by whether or not these transitions require phonon participation. Direct transitions (i.e. transitions between states with the same wavevector) can occur without phonon participation, and can have relatively large transition strengths. However, indirect transitions (i.e. transitions between states with unequal wave vectors) require phonon emission or phonon absorption in addition to the photon emmision or absorption process, and have relatively weak transition strengths.\nThe efficiency of optical devices can be significantly improved by the use of active region materials with direct transitions at the photon energies of interest. In addition, certain optical devices such as laser diodes essentially require the use of materials with direct band-edge transitions; i.e., electron-hole recombination events which occur at energies close to the band gap.\nMany materials having direct transitions are used in optical devices. For instance, bulk GaAs and InAs are examples of conventional direct transition materials. However, some applications require the use of materials having indirect bandgaps. For example, all of the known technologically relevant compound semiconductors with appropriate bandgaps for green LEDs and green laser diodes have indirect band gaps. The emission efficiencies of sub-bandgap transitions can be enhanced somewhat in certain materials by the incorporation of selected impurities. However, the efficiencies obtained in this manner do not approach those obtained from direct gap materials."}
-{"text": "This invention relates to a method of transmitting information between a plurality of stations in a local area network.\nVarious methods are known for sharing network communication channels among stations in a local area network (LAN). One widely used method is known as CSMA/CD, (carrier sense multiple access/collision detect). According to this known method, a station wishing to transmit a message listens until the transmission channel is idle before commencing to transmit an information packet. Furthermore, the station continues to listen to the channel after commencing a transmission and if a collision is detected, that is, more than one station has commenced to transmit an information packet, any station which detects such collision terminates its message transmission and transmits a burst of noise (or garble) such that all other stations can be informed of the collision, terminate message transmission, and wait a random time before attempting to commence a further message transmission.\nFrom Patent Abstracts of Japan volume E-690, page 102, abstract of Japanese Published Application No. 63-187746, there is known a wireless communication system wherein a plurality of stations communicate via a repeater installation. A first frequency is used for transmission from the stations to the repeater installation and a second frequency is used for transmissions from the repeater installation to the stations. If the repeater installation detects a collision on the first frequency transmissions, then it immediately informs all stations of such collision using the second frequency, such that the stations can take appropriate action.\nAlthough the CSMA/CD method has proved satisfactory for transmissions over a cable physical medium, it has been found that a problem arises if it is attempted to apply the method to wireless radio transmissions in that small radio signals may be totally masked by larger signals during reception. This phenomenon is known as the \"capture effect\", and may prevent the detection of collisions in a system such as that disclosed in the aforementioned Japanese patent abstract."}
-{"text": "1. Field of the Invention\nThe presently disclosed subject matter relates to a technique that retrieves and presents a case image similar to a diagnostic object image (image of object to be diagnosed) on the basis of an amount of feature (feature amount) extracted from the diagnostic object image.\n2. Description of the Related Art\nIn Japanese Patent Application Laid-Open No. 2004-5364, an ROI (region of interest) is specified on image data; information such as image data which has a similar feature and a diagnostic result of the similar image data are retrieved from a database and displayed.\nIn Japanese Patent Application Laid-Open No. 2001-117936, a region of interest is extracted from a diagnostic image, and a similarity search is performed using the amount of feature of the region of interest.\nJapanese Patent Application Laid-Open No. 2008-257292 shows another example of similar image search system.\nJapanese Patent Application Laid-Open No. 2008-245719 shows an example of a technique for extracting a lesion region on the basis of position information of a region of interest.\nT. F. Cootes, G. J. Edwards, and C. J. Taylor \u201cActive Appearance Models\u201d, In Proc. 5th European Conference on Computer Vision, Springer Germany, volume II: pp. 484-498 (1998), shows a method of acquiring shape information and texture information of a lesion part in the ROI (region of interest)."}
-{"text": "The present invention relates to Caterpillar sidebooms used for pipelaying, and more particularly pertains to methods and apparatus for converting old, discontinued, mechanically-operated, low-drive, Caterpillar sidebooms into hydraulically-operated sidebooms, with minimal retrofitting.\nIt is well known in the art that crawler or tractor-type vehicles having an integrated, rigid maneuverable boom disposed on a side thereof are commonly used for pipelaying operations, i.e., for raising, carrying, and lowering heavy pipe. Referred to as \"sidebooms,\" such vehicles must be capable of safely handling heavy pipe; indeed, there are sideboom models that have been constructed for handling pipe up to 200,000 pounds. Such sidebooms were mechanically-operated in their first incarnation and are now readily available from original equipment manufacturers in hydraulically-operated models.\nAn early development in the mechanical sideboom art is disclosed by Butterfield et al. in U.S. Pat. No. 3,785,503, wherein a plurality of planetary gears and concomitant shafts are used to drive the winch drums for each of the boom and the load. Primarily due to its inherent complicated, high-maintenance gearing system, no embodiments of the Butterfield sideboom was ever commercialized. A significant improvement in the pipelaying sideboom art would be a sideboom that eliminates all gears to accomplish transmission of power, i.e., that eliminates all mechanical connections between the engine and the winch system.\nOther early developments in the sideboom art, albeit not applicable to pipelaying operations, per se, are taught in U.S. Pat. Nos. 2,909,290; and 3,329,283. More particularly, Nichols, in U.S. Pat. No. 2,909,290, teaches a farm tractor-mounted sideboom intended for lifting light loads typical on the farm. Of course, pipelaying operations demand sidebooms with lifting capacities up to 200,000 pounds. Similarly, Wade, in U.S. Pat. No. 3,329,283, teaches a snap mount sideboom configured to be foldable for reducing its prerequisite overhead clearance. Half of the Wade boom's height may be reduced and the folded boom portion secured to the tractor's side by using integral hooks and cable. Especially in view of the Wade sideboom having no counterweight, it should be evident to those skilled in the art that both Wade and Nichols are inapplicable to the rigors of pipelaying wherein not only a strong, firm boom structure is required, but also the framework of the sideboom must be sufficiently broad to provide a low enough center of gravity for stability and must include a counterweight to provide sufficient operational stability and safety. It will be readily understood by those skilled in the art that such attributes are not provided by conventional tractors.\nIndicative of initial attempts to improve the sideboom art using hydraulics is an apparatus described by Stefanutti in U.S. Pat. No. 3,265,218. In particular, the use of hydraulically actuated booms and hoist assemblies is described, wherein hydraulic cylinders are used either inside or astride the boom to raise and lower the boom. As will be appreciated by those skilled in the art, the Stefanutti apparatus has not been widely accepted.\nAs taught by Vinton in U.S. Pat. No. 3,938,669, however, while the introduction of such hydraulically-operated sidebooms were anticipated to provide improvements associated with effectively and safely manipulating winches and cables through an operator's interfacing with clutches, brakes, and levers for controlling the position of the boom and for hoisting and lowering a load, such sidebooms failed to provide the prerequisite control and versatility. To attempt to remedy this deficiency in the art, Vinton discloses a hydraulic circuit that includes two separate sources of hydraulic fluid for controlling sideboom movement-related functions. One fluid source provides low volume hydraulic fluid for accomplishing not only precise, low-speed manipulation and control, but also for preventing anti-drift of both boom and hoist. The other fluid source provides high volume hydraulic fluid for providing high-speed operation of these movement and control functions. The plurality of control valves inherent in the Vinton circuit for controlling the hoist and the boom motors are operated via two levers. This apparatus also incorporates a hydraulic cylinder, instead of a drum/cable arrangement for controlling the boom; no drums or cables are involved in raising or lowering either the boom or the hook, except a short cable connected to the hook at the end of a hydraulic cylinder. As is well known by those skilled in the sideboom art, this methodology was a commercial failure.\nAs a further development in the sideboom art, Forsyth teaches in U.S. Pat. No. 5,332,110 a hydraulically-operated sideboom intended to prevent boom over-rotation, to impart positive drive to the boom and load winches, and to provide improved control over free fall and vertical kick-out. These safe operating features are particularly intended for pipelaying applications involving lifting and lowering of large pipes. Indicative of current Caterpillar tractors and sidebooms, embodiments taught by Forsyth are exclusively for newly manufactured hydraulically-operated high-drive tractors which are the antithesis of predecessor low-drive mechanically-operated Caterpillar tractors.\nAs will be appreciated by those skilled in the sideboom art, current Caterpillar hydraulically-operated high-drive sidebooms have a higher center of gravity but operate easier due to less controls than old, discontinued, mechanically-operated, low-drive, Caterpillar sidebooms. Such \"high drive\" models are constructed with an undercarriage that was introduced by Caterpillar in the early 1980's; the sprockets and drive mechanisms for the undercarriage are situated much higher off the ground than the previous conventional \"low drive\" system that was in effect from the inception of the original Caterpillar sideboom. More particularly, the Caterpillar low-drive sidebooms were discontinued from 1972 through 1986 as follows: model 572E was discontinued in 1972; model 572F was discontinued in 1975; model 572G was discontinued in 1986; model 583H was discontinued in 1974; model 583K was discontinued in 1986; model 594G was discontinued in 1975; and model 594H was discontinued in 1986. In addition, Caterpillar hydraulically-operated high-drive sidebooms are very expensive, and some models not only have inherent counterweight obstruction problems, but also are difficult to move from job-site to job-site. While the Forsyth disclosure teaches that his hydraulic pipelayer is adapted for mounting upon a conventional track-laying tractor (i.e., bulldozer), it is well known in the art that the main frame of a sideboom is constructed differently from that of a conventional tractor. In particular, unlike a conventional tractor which is constructed with an oscillating frame, a sideboom is constructed with a rigid frame of wider track gauge than a conventional tractor. Indeed, Caterpillar identifies such tractor and sideboom frames with different serial numbers series. Thus, to obtain the prerequisite performance demanded in the pipelaying art, a drawworks assembly must be mounted upon a frame capable of rigidity to accommodate the pivoting action of a sideboom typically positioned upon rough terrain, with the frame having a sufficiently wide track gauge for stability purposes.\nAs will be understood by practitioners in the art, a drawworks system built upon an old, discontinued, mechanically-operated, low-drive, Caterpillar sideboom--having significantly more controls than a conventional discontinued tractor--inherently suffers from a panoply of problems associated with the simultaneous use of a daunting ensemble of gear-shifting mechanisms, clutches, and brakes, all operated by 6 different hand controls to properly lift and manipulate heavy pipes under conditions generally characterized by unpredictable and adverse terrain. Pipeline construction companies constitute 95% of the users of this type of machinery. As will be appreciated by those conversant with the art, historically, such pipeline construction companies have had to choose between the newer high-drive sidebooms with herein before mentioned faults and high price or the older, discontinued Caterpillar sidebooms that are more economical but are more dangerous and are very difficult to operate. It should also be noted that, as the pipelaying industry continues to mature, the number of skilled sideboom operators has gradually diminished.\nAs should be evident to those skilled in the art, it would be advantageous for construction companies who are continuing to utilize older discontinued Caterpillar sidebooms to have the additional benefits of improved handling, safety, and efficiency. It would be also be advantageous for pipeline contractors to have the ability to expeditiously train low-drive Caterpillar sideboom operators and to simultaneously achieve a level of safety heretofore unmatched by any other sideboom system known in the art, regardless of design. Of course, it would be advantageous for pipelaying contractors to have the benefit of a Caterpillar sideboom that inherently avoids or mitigates the complex levers and the like associated with maneuvering a sideboom, and controlling the lifting and lowering of a pipe load.\nAccordingly, these limitations and disadvantages of the prior art are overcome with the present invention, and improved means and techniques are provided which are useful for effectively and reliably utilizing old, discontinued, mechanically-operated, low-drive, Caterpillar sidebooms for pipelaying applications."}
-{"text": "1. Field of the Invention\nThe present invention relates to apparatus and method for verifying an integrated circuit pattern formed on a substrate.\n2. Description of the Related Art\nIn the manufacture of semiconductor devices, there is a need of verifying whether an integrated circuit pattern formed on a substrate is formed as designed in individual processes.\nFIG. 15 is a view to explain a method of verifying a hole pattern. According to the method, a target data TD of preset hole pattern is overlapped with a contour data CD extracted from hole pattern formed on a substrate. In the overlapped state, distances from several places of the foregoing two patterns are measured. It is determined whether or not the difference between the measured distances is within a predetermined allowable value, and thereby, it is determined whether the formed pattern is non-defective or defective.\nThe method of evaluating a pattern shape has been proposed (for example, JPN. PAT. APPLN. KOKAI Publication No. 2002-31525). According to the method, the distance between CAD (Computer Aided Design) output data and pattern data measured by SEM (Scanning Electron Microscope) is measured in the same manner as described above. Based on the measured distance, evaluation is made with respect to the pattern shape.\nThe method of evaluating a pattern defect information has been proposed (for example, JPN. PAT. APPLN. KOKAI Publication No. 2002-93875). According to the method, the error dimension between CAD output data and shape simulation output data is measured. If the measured error dimension is not within a predetermined value, pattern defect information is displayed.\nIn addition, the following method has been proposed (for example, JPN. PAT. APPLN. KOKAI Publication No. 2000-182912). According to the method, comparison is made between an inspection reference pattern formed based on the design layout pattern and the finished predictive pattern. By doing so, pattern distortion of the finished predictive pattern is detected, and thereafter, the detected pattern distortion is identified in accordance with criticality.\nAccording to the method of measuring the distance between the foregoing two patterns, a great many of places is measured, and thereby, the measurement accuracy is improved. However, if the measured places are too many, long time is spent for measuring these places. According to the method, the difference from the corner of the pattern is not so taken into consideration. For this reason, when verifying a pattern having a complicated shape, it is very difficult to determine whether the pattern is non-defective or defective. Moreover, there exist various patterns requiring to bridge neighboring patterns depending on pattern, or to control only line width of the shorter side of the pattern. In the foregoing patterns, the allowable (variable) range for the target dimension has not been clarified. For this reason, according to the conventional method, it is determined that patterns formed with accuracy having no problem on the device operation are defective.\nAccording to the conventional method, time is taken to verify the pattern, and the features of the pattern are not considered; as a result, the pattern is not properly verified. In addition, according to the conventional method, it is difficult to evaluate a complicated pattern."}
-{"text": "The concept of the \u201cintelligent network\u201d (IN) was developed for efficient deployment of resources when using digital communication networks. This is a structural concept for flexible service control and service development through the use of software and functionality implemented therewith. Internationally IN is defined in the Q.1200 series recommendations, for which the ITU (International Telecommunications Union) is responsible. The basic thinking behind IN is the separation of switching functions and service functions, i.e. resources for connection set-up are no longer located wholly in the individual switching point but are centralized. The advantage of this concept is simpler management and adaptation of these centralized resources. For example, when new services are introduced, not all switching points have to be adapted individually.\nThere are three hierarchical levels in the IN architecture, management level, service control level, and switching and transport level. The switching points\u2014generally referred to as SSP (service switching point)\u2014of the lowest level, i.e. the switching and transport level, are controlled by control points\u2014generally referred to as SCP (service control point)\u2014at service control level for the implementation of services and where appropriate are supplied with data from centralized databases. SSP and SCP are defined as physical entities for the IN concept in Q.1205.\nGeneral aspects of an INAP (intelligent network application protocol) protocol are described in Q.1208 for IN applications. At present INAP protocols with CS2 (capability set 2) and CS3 (capability set 3) operational features in accordance with the ITU recommendations numbered Q.122x and Q.123x are frequently used for IN applications.\nFurther specifications are required to provide telecommunication services and service features in the context of the IN concept. For example, for ISDN (integrated services digital network) applications signaling functions are standardized in the context of the ISUP (ISDN user part) in ITU documents Q.761-Q.764 and Q.766, with the ISUP representing a protocol based on signaling system no. 7 (SS#7), which allows separation of the transfer of user data and control data and which is standardized by the ITU in the Q.700 series. For ISDN applications in the context of the IN concept the ITU has produced standards Q.1600 and Q.1601 which govern the interaction of ISUP and INAP CS1 or INAP CS2.\nGiven the increasing importance of the transfer of user data and any signaling data via different networks or carriers, as is frequently the case, for example, with IP based applications, signaling protocols are being developed, which allow a specific set of services and service features, regardless of the carrier. An important example of such a service is internet telephony (VoIP), in which connection-based voice transfer is to be achieved via different physical networks. The ITU draft of Q.1901 describes a signaling protocol, the BICC (bearer independent call control) protocol, with which the services and service features defined for the ISDN network\u2014hereafter referred to as ISDN services or service features\u2014are to be implemented, regardless of the carrier. The BICC protocol is based on the ISUP protocol. As described in the draft of Q.765.5, the transport mechanism (application transport mechanism) set out in Q.765 can be used for the transport of carrier-related data using the BICC protocol.\nWhen implementing ISDN services and service features in the context of the IN concept with user data transfer via a packet network using a carrier-independent signaling protocol, such as, for example, BICC CS2, there is the problem that the parameter changes necessary or desirable for this are not supported."}
-{"text": "Polishing is an activity which is useful to maintain the appearance and condition of footwear, and other accessories such as belts, purses, and the like. These objects, whether made from fabric, leather, or similar materials, often become dirty or dull as a result of use. Various implements have been developed for applying treatment agents to the surfaces of such objects to clean or enhance their finish, and for polishing after treatment. In many cases, more than one implement is required to complete the polishing process; for example, when polishing shoes, it is often necessary to use one implement for applying polish to the shoes, and another implement to buff the polish in order to achieve the desired shine. Moreover, such implements generally require both hands, and use of them often results in transfer of treatment agent to the hands of the user. Because of the need to use more than one implement, and the mess that is often associated with the process, many users are not inclined to regularly polish their shoes. Accordingly, there is a need for a shoe care apparatus which is multifunctional, can be used with one hand to achieve effective treatment and buffing, and which minimizes the need for the user to come in direct contact with polish and other treatment agents."}
-{"text": "Memory technology implies the need for accessing data within memory devices, such as DRAM or other memory devices. Therefore, different addressing strategies have already been proposed.\nA possible feature of DRAMs is address multiplexing. This technique enables splitting the address in half and feeding each half in turn to the chip on the address bus pins.\nThe chip has a large array of memory capacitors that are arranged in rows and columns. To read one location in the array, the control circuit first calculates its row number, which it places on the DRAM's address pins. It then toggles the row address select (RAS) pin, causing the DRAM to read the row address. Internally, the DRAM connects the selected row to a bank of amplifiers called sense amplifiers, which read the contents of all the capacitors in the row. The control circuit then places the column number of the desired location on the same address pins, and toggles the column address select (CAS) pin, causing the DRAM to read the column address. The DRAM uses this to select the output of the sense amplifier corresponding to the selected column. After a delay called the CAS access time, this output is presented to the outside world on the DRAM's data I/O pin.\nTo write data to the DRAM, the control logic uses the same two-step addressing method, but instead of reading the data from the chip at the end of the operation, it provides data to the chip at the start of the operation.\nAfter a read or write operation, the control circuit returns the RAS and CAS pins to their original states to ready the DRAM for its next operation. The DRAM requires a certain interval called the precharge interval between operations.\nOnce the control circuit has selected a particular row, it can select several columns in succession by placing different column addresses on the address pins, toggling CAS each time, while the DRAM keeps the same row activated. This is quicker than accessing each location using the full row-column procedure. This method is useful for retrieving microprocessor instructions, which tend to be stored at successive addresses in memory.\nIn addition, the provision of commands is typically provided through command strobes on the command bus. The number of different commands depends on the number of pins on the command bus with 2N commands being a possibility with N being the number of pins at the command bus. With the increasing demand for different commands, the command bus needed to be expanded. However, as die size is a crucial factor in application specific integrated circuit (ASIC) design, the number of pins on the command bus needs to be decreased.\nFurther, the size of the buses is also relevant for the overall size of the connection interface between the memory device and the central processing unit (CPU). The higher data rates that were required, the higher the number of connection pins on the data bus were selected. This increased the size of the interface. In addition, the number of pins on the address bus and the overall number of pins of the interface determined the type of memory to be used on the interface, besides protocol issues. However, the demand for flexibility of usage of different kinds of memory devices was not accounted for. There is a need for a flexible interface, which enables the use of different kinds of memory devices with different kinds of capabilities in terms of data throughput on the data bus.\nWith the rising need for flexible use of standard components, there is a need for providing a memory interface enabling to use both volatile memory and non-volatile memory on one same interface. However, as non-volatile memory and volatile memory have different prerequisites for the interfaces, there needs to be a possibility to adapt the interface for use with both types of memory."}
-{"text": "The present invention relates to hair thickening and conditioning compositions, and in particular to aqueous gelatin solutions, the application of which to the hair results in a significant increase in hair shaft thickness. The present invention also relates to methods for preparing the gelatin compositions of the present invention, as well as to methods for increasing hair shaft thickness by topical application to the hair of the gelatin compositions of the present invention.\nU.S. Pat. No. 4,749,684 discloses that the ingestion of gelatin promotes an increase in the linear growth rate of hair. Scala et al., Nutr. Rep. Int., 13(6), 579-92 (1976) discloses that dietary supplementation with gelatin produces an increase in hair thickness. The ability to increase hair thickness by the topical application of gelatin is unreported.\nHair coated with gelatin is difficult to comb in the absence of a conditioning agent, but stable solutions of gelatin with a conditioning agent have been heretofore unknown. Lindo et al., \"Hair Processing and Conditioning,\" Hum. Hair Sym. [Pap.], 1st, (Brown, Editor, MEDCOM Press, New York 1974 (Meeting Date 1973)) pages 135-44 reports the capability of hair to absorb topically applied hydrolyzed gelatin, but with little substantivity, so that little gelatin remains after thorough washing with room temperature water.\nU.S. Pat. No. 4,208,402 discloses a topical solution for improving the appearance of oily hair or skin, including the scalp, containing from about 0.5 to 1.5% by weight of gelatin and from 0.15 to 0.45% by weight of partially hydrolyzed polyvinyl acetate in aqueous or hydroalcoholic solution. The level of film-forming polymer disclosed is insufficient to provide a readily combable product. By the manufacturing process disclosed, aqueous solutions cannot be formed with higher concentrations of the film-forming polymer without the gelatin precipitating from the solution.\nTopical application of gelatin to hair would be possible if stable aqueous solutions containing higher concentrations of film-forming polymers could be prepared. The inability to comb hair to which gelatin solutions have been topically applied has prevented the recognition that topical gelatin application increases hair shaft thickness.\nThere remains a need for stable aqueous gelatin solutions containing higher levels of film-forming polymer to provide gelatin solutions for topical hair application that do not resist combing."}
-{"text": "1. Field of the Invention\nThe present invention generally relates to high voltage on-chip circuitry for electrically erasable programmable read-only memory EEPROMs or flash memory. More particularly a low power, high-voltage ramp-up rate control circuit is described which utilizes a current reference, a current controlled oscillator, a charge pump, a feedback transistor, a reference transistor, and a ramp-up transistor.\n2. Description of the Prior Art\nFIG. 1 contains a circuit schematic which shows the prior art implementation of a ramp-up rate control circuit. A functional block labeled charge pump+oscillator+current reference 110 is shown. The output of this block 110 is a charge pump current, Icp 170. In addition, the voltage on the output node 195 of this block 110 is HV. (High Voltage). A reference current, Iref, 120 is developed in a current source. This current flows into the drain of an N-channel metal oxide semiconductor field effect transistor, NMOS FET, M1130. The gate 175 of FET M1130 is connected to its drain. The source of FET M1 is connected to ground 185.\nThe ramp current Ir 190 flows into the drain 165 of NMOS FET M2140. The gate of FET M2140 is connected in common to the gate of FET M1175. The source of FET M2 is connected to ground 185. The drain 165 of FET M2 is connected to the gate of NMOS FET M3150. The drain 180 of FET M3150 receives the leaking current Ileak. The source of FET M3 is tied to ground 185.\nIn this prior art circuit in FIG. 1, the leak transistor M3150 will tend to discharge the charge pump output node HV 195. The leakage current Ileak 180 depends on the gate voltage V1165 of FET M3150. The ramp-up rate (dV/dt) is equal to Ir/C where C is the capacitance of the capacitor 160 and V is the voltage across capacitor C. In FIG. 1, the leakage path will consume the extra current Ileak. If this extra leakage path 180 can be eliminated, then the power consumption can also be reduced.\nU.S. Pat. No. 4,488,060 (Simko) \u201cHigh Voltage Ramp Rate Control Systems\u201d describes high voltage ramp rate control systems. These systems are useful for on-chip EEPROM high voltage power supplies.\nU.S. Pat. No. 5,872,733 (Buti et al.) \u201cRamp-Up Rate Control Circuit for Flash Memory Charge Pump\u201d describes a ramp-up rate control circuit. This circuit controls the ramp-up rate of a charge pump having an output which provides an output voltage and an output current.\nU.S. Pat. No. 5,408,133 (Honnigford et al.) \u201cRamp Control Circuit\u201d discloses an apparatus for providing an EEPROM programming signal.\nU.S. Pat. No. 5,365,121 (Morton et al.) \u201cCharge Pump with Controlled Ramp Rate\u201d discloses a charge pump with a controlled ramp rate. It is made up of a charge pump, a differentiator circuit and a trigger circuit."}
-{"text": "Many satellites and other systems have one or more covers or doors to protect internal mechanisms, e.g., optics, sensors and sensitive machinery during transport or other periods of disuse. These types of doors utilize various mechanisms, such as mechanical launch locks to hold them in a closed position and motors to open them. Motors and the control circuitry by which they are operated are expensive to design, procure, manufacture and test. In some applications, such as space-based systems, testing typically includes a life cycle test which could take months to complete. Additional items that drive cost in such space-based systems include the launch lock, the control/telemetry electronics, the drive motor, the gear box, as well as vacuum compatible bearings and grease. Research and development efforts that seek to simplify cover/door mechanisms, and to reduce cost accordingly, are ongoing."}
-{"text": "The present invention relates to a method for analyzing a stereoscopic image of a fundus to detect and analyze the optic disk on the basis of the stereoscopic image of the fundus, and an apparatus for executing that method.\nGlaucoma is known as a disease of the visual function. In the examination of glaucoma, it is said that the observation of patient's optic disk on the basis of an image of the fundus is important. In recent years, it has become possible to quantitatively measure the optic disk due to progress in the technology of image analysis. As the results of analysis, the following parameters are computed: the C/D ratio (a distance ratio between a rim of an optic disk cupping (the optic disk cupping being referred to as a cup) and a rim of the optic disk (the optic disk being referred to as a disk)), the area of the plane within the cup rim, the depth of the cup, the volume of the cup, and the like.\nThe rim of the cup which serves as a basis for the computation of these parameters is conventionally detected such that, by using the disk as a reference, positions which are lowered by an arbitrary depth (e.g., 150 .mu.m) from the disk are estimated to plot the rim of the cup.\nHowever, since it is not uncommon that overall disk regions of the patent's eyes are inclined, there arise cases where the rim of the cup cannot be detected by the above-described detection method.\nIt is possible to obviate the above-described problem if an operator designates three points on the disk region and effects analysis by using as a reference plane an inclined plane thus determined. In accordance with this method, however, there is a problem in that the data of the cups detected vary for each operator, and the data can be different on each occasion of analysis even for the same operator, with the result that the reliability and reproducibility of the results of analysis are poor."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to a self contained lamp control and lamp package, and relates more specifically to a self contained lamp for automatic light level adjustment.\n2. Description of Related Art\nReferring to FIG. 1, a conventional incandescent light bulb 5 includes base 10, and evacuated envelope or diffuser 14. Diffuser 14, which is typically pear-shaped, surrounds a filament (not shown) that is electrically connected to base 10 for electrical connection to a conventional lamp socket. A well-known base is an Edison screw-base, as shown, which is electrically connectable to an Edison screw socket.\nIn recent years, compact fluorescent lamp packages have been introduced into the market. These lamp packages are particularly desirable because of their energy efficiency.\nReferring to FIG. 2A, an example prior art compact fluorescent lamp package (CFLP) 7 is shown consisting of fluorescent lamp 12, lamp holder 13, and auxiliary housing 15 for interfacing between lamp holder 13 and base 10. Typically, housing 15 also houses a complete ballast 17. However, it should be noted that ballast 17 has also been placed into base 10, while either retaining housing 15 or removing this housing. In most CFLPs an Edison screw-base is used as base 10 for a reliable mechanical and electrical connection to a standard Edison socket, although other sockets are known and used.\nThe typical height of a standard incandescent light bulb 5 is about 4.5 in. However, most conventional CFLPs are typically taller than a conventional incandescent light bulb due to auxiliary elements such as auxiliary housing 15. As a result, a conventional CFLP may extend out of a lampshade or fixture adapted for an incandescent light bulb, causing an undesirable appearance, or the CFLP may not fit within an existing fixture. This extended height, together with the unusual appearance of conventional CFLPs, are believed to adversely affect the acceptance of CFLPs in the marketplace\nAnother issue with the conventional CFLP is that when ballast 17 is housed in auxiliary housing 15, the heat generated by ballast 17 cannot escape very efficiently. Referring, for example, to FIG. 2B, in a typical enclosure or downlight fixture 19, the heat generated from the lamp and ballast of a conventional CFLP 7 collects up in the fixture around the CFLP 7 without any path to escape.\nAs a result, there is a higher ambient temperature around CFLP 7 which causes the internal components of the CFLP 7, especially ballast 17, to run at even higher temperatures, e.g., above 150\u00b0 C. The increase in the internal heat decreases reliability, causes field failures and limits the use of a conventional CFLP to open rather than enclosed fixtures, all of which further adversely affects the desirability of CFLPs. Nonetheless, it should be noted that when ballast 17 has been placed into base 10 as described above, the base has acted as a heat sink for the ballast, helping to dissipate some of this heat.\nA further issue with the conventional CFLP is that because of its size and operational constraints, such as its high ambient temperature and its use of the conventional Edison socket, it is difficult to realize some desirable features, such as, for example, a dimming feature where the light output of the fluorescent lamp is selected based on a control signal or setting. More specifically, conventional incandescent lights typically use a dimmer control mounted on a wall near an on/off switch for adjusting the light level. These types of dimming controls operate on various principles but all in essence change the characteristic of the electrical power supplied to the light fixture to create various power levels supplied to the lamp, thereby affecting the output lighting level. However, these types of dimming controls are unsuitable for CFLPs because the change in power may extinguish the lamp. As a result, in order to incorporate a dimming feature into the CFLP, a dimmable ballast circuit is typically used where the electronic ballast now includes an additional dimming function that controls the ballast to dim the lamp. However, a mechanism is now needed to control this additional dimming function, thereby making it a difficult challenge to provide a dimming feature to a CFLP.\nFor example, one mechanism that has been used to control a dimmable ballast circuit is to dedicate a wire from a wall-mounted dimmer to the fixture. However, if the CFLP is inserted into an Edison socket, this requires additional wiring. Specifically, Edison screw sockets typically only provide switched utility power on two wires, a line input and a neutral wire. Accordingly, a third wire is now needed to control the dimming function. However, providing a third wire or connection to the Edison screw socket is a highly undesirable additional requirement to realize a dimming feature.\nAnother example mechanism that has been used to control a dimmable ballast circuit is to send control signals from a switch to the fixture directly over the power lines, thereby negating the need for an additional wire. However, this approach is complex, expensive, and space consuming given that the fixture must now include an interface to receive and decode the signal. Nonetheless, it should be noted that some systems that have used this mechanism have also incorporated advantageous features, including the use of an external control module that allows for the remote and automatic control of the lights. For example, a control system may be provided that can control a number of different areas of lighting to automatically change lighting settings depending upon various circumstances, such as evening hours, weekend hours, lengthened days and so forth. However, it should also be noted that these types of lighting systems typically have limited local intelligence (e.g., the fixtures lack sensor feedback that permits localized control) and often control blocks of lights together, rather than individual lights.\nAnother mechanism that has been used to control a dimmable ballast circuit is to include a manual control directly on the casing of the electronic ballast of the CFLP so that dimming can be effected. However, CFLPs designed to replace incandescent bulbs often do not permit easy access to the CFLP or electronic ballast by the user. In addition, once the light level has been set, it is often difficult or inconvenient to readjust.\nA further mechanism that prior fluorescent light systems in general have used to control dimming is to incorporate a self-dimming function that automatically dims the light in reaction to external events. Such self-dimming functions have included the integration of a heat sensor or a photosensitive control element, such as a photocell, with the dimming function. For example, these systems have interconnected a photocell with the dimming function to automatically control the turning on/off of the light and to maintain a constant light level in a room, all based on the changing ambient light levels in the room.\nIn the case of the CFLP in particular, the photocell has been mounted directly on the housing 15. This arrangement creates several advantages as compared to the other dimming mechanisms described above. Specifically, this arrangement allows for individual control of each CFLP and because of the automatic feedback orientation, each CFLP is self-dimming thereby making the CFLP easier and more convenient to use because the user is not required to constantly readjust the light level as conditions change. In addition, when the CFLP is inserted into an Edison socket, the socket does not require additional wiring. This mechanism is also less complex and less expensive than prior mechanisms.\nNonetheless, prior CFLPs that have integrated a sensor, like a photocell, with the dimming function of a dimmable ballast circuit continue to have the same issues as described above with reference to the example CFLP 7. Specifically, these prior CFLPs continue to use a housing 15 to house the dimmable ballast circuit, thereby creating the same height and heat issues described above."}
-{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display (LCD), and more particularly, to a response time accelerator system and method for driving a liquid crystal display.\n2. Description of the Related Art\nCurrent LCD technologies have problems associated with liquid crystal response time. Thus, because the response time of the liquid crystal forming pixels in an LCD panel is relatively slow, a user sees an after image, when a TV displays a large number of moving images.\nEach liquid crystal cell within an LCD panel allows light to pass through or blocks the light when an applied bias voltage rotates a lattice. Usually, liquid crystals cannot respond in real-time to data changes since the time required to respond to the bias voltage is on the order of tens of milliseconds and bias voltage applied across the crystal varies from frame to frame (1/75 seconds at SXGA resolution). For example, if the bias voltage applied across the liquid crystal within an LCD panel is set to a gray level of 255 for 8-bit image data, the brightness level after the actual response of the liquid crystal is less than 255, which causes vertical stripe patterns to generate an after image. The liquid crystal response characteristics is aggravated as the frame period decreases due to increased resolution.\nThe most commonly used method for preventing liquid crystal panel after image retention due to a slow response is to appropriately preprocess image data prior to processing it in a source driver for driving a liquid crystal panel. To implement this technique, a response time accelerator (RTA) has been adopted. However, the RTA has many problems.\nIn order for the liquid crystal to keep up with current frame data, the RTA basically compares a previous frame data with the current frame data, interpolates the two frame data and finds a new current frame data with respect to which a response time can be accelerated according to the result of the comparison. The current data input to the RTA in real time is compared with the previous frame data stored in a memory such as SDRAM outside the RTA.\nA method for accelerating a response time applied to a conventional RTA involves storing four to six most significant bits (MSB's) of RGB (Red, Green, Blue) data as frame data in an external memory such as SDRAM. However, when storing only the MSBs (e.g., n MSBs) in an external memory such as SDRAM as frame data, a pixel data truncation error (PDTE) occurs. If only the MSBs are used as previous frame data Pn\u22121, the number of gray level levels that can be compared between the previous frame data Pn\u22121 and current frame data Pn is (2n\u00d72n) where n bits Pn\u22121[7:8\u2212n] of previous frame data Pn\u22121 and n bits Pn[7:8\u2212n] of current frame data Pn are used. This causes an error to occur. That is, a quantization error occurs because the number of gray levels available for comparison is reduced to (2n\u00d72n) from \u201c256\u00d7256\u201d, which is the number of gray levels available when 8-bit data is entirely used for each frame data. Here, the term \u201cquantization error\u201d is given because the error occurs sporadically due to unused 8-n LSBs.\nFor example, assuming a value derived by shifting n bits Pn[7:8-n] of the current frame data Pn to the right by 8-n bit positions is K (e.g., Pn[3:0] if n=4), the quantization error occurs when outputting an overshoot value exceeding or an undershoot value below the current frame data to the panel at gray level K*2(8\u2212N). Eventually, since the quantization error occurs periodically at gray level produced by multiplying 2(8\u2212N) by the value derived by shifting current frame data for each gray level by 8-n bit positions, noise is generated at regular intervals in vertical stripes to be displayed on a liquid crystal panel. The above description, as previously mentioned, is under the assumption that the LCD uses 8-bit RGB data for each pixel making it possible to display 256 gray levels.\nHere, the new current frame data found by interpolation of the previous frame data Pn\u22121 and the current frame data Pn is usually stored in a table memory within an RTA. The tabular values stored in the built-in table memory for interpolation, defining overshoot exceeding or undershoot below the current frame data, are experimentally determined based on the liquid crystal characteristics of a panel. For example, given the fact that the response time of liquid crystal in the panel is relatively slow, if the current frame data is greater than the previous frame data, the new current frame data is assigned a tabular value greater than actual data, which is then output to the panel. Similarly, if the current frame data is less than the previous frame data, the current frame data is assigned a tabular value less than the actual data, which is then output to the panel.\nAnother drawback of the method for accelerating a response time implemented using a conventional RTA is that it limits the amount of overshoot and undershoot for each output data to a gray level range from 0 to 255 since the number of bits of RGB image data is limited to 8. Therefore, if the current frame data Pn has an extremely large value (around a gray level of 255) or an extremely small value (around a gray level of 0), the output data does not have sufficient amount of overshoot or undershoot. For example, if the current frame data Pn has a maximum gray level of 255, it is impossible to give an overshoot greater than the current frame data Pn to the output data since the overshoot is limited to the maximum gray level of 255. Thus, if data output from the panel is at a gray level of 255, the liquid crystal panel will return a value less than 255 in response. Such an overshoot or undershoot limitation makes it difficult to improve a response time."}
-{"text": "1. Field of the Invention\nThe present invention relates to a drive control technique for an optical modulator used in optical communication. In particular, the invention relates to a drive control apparatus and drive control method, which can stably control the drive of an optical modulator used in ultra high speed optical transmission and the like.\n2. Description of the Related Art\nRecently, in optical transmission systems, with the progress of wavelength division multiplexing (WDM) technology and higher transmission speed per wavelength, systems enabling of large capacity optical transmission are being realized. Also, for the modulation code form for transmission used in such systems, in order to enlarge the transmission capacity and transmission distance, forms such as optical duo binary code and carrier suppressed RZ (CS-RZ) code have been proposed, instead of conventional NRZ (non return to zero) code or RZ (return to zero) code (refer to Japanese Unexamined Patent Publication No. 2001-119344). In particular, the CS-RZ code is said to be the most promising modulation code form for the next generation, since it enables enlargement of transmission distance and transmission capacity by a simple method.\nAlso, in optical modulation for ultra high speed transmission, in order to suppress wavelength fluctuations (chirping) during modulation, an external modulation method has been used, for allowing a semiconductor laser to continuously emit light, to switch the light on and off by an external modulator. For the external modulator, a Mach-Zehnder optical modulator (to be referred to as MZ modulator hereunder) is prevailing.\nFIG. 31 shows an example of optical output characteristics of an MZ optical modulator with respect to a drive signal. In the conventional NRZ modulation method, on/off modulation of the optical output is performed by matching a high level and a low level of a drive signal, whose level fluctuates at a period T0, with an apex A of optical emission and an apex B of optical extinction in the optical output characteristics. In the following explanation, a difference of bias voltages corresponding to the apexes A and B of emission/extinction in the optical output characteristics which fluctuate periodically is represented by V\u03c0. According to this, an amplitude of the drive signal in the NRZ modulation method described above becomes V\u03c0.\nFor the MZ optical modulator as described above, the advantage is less chirping. However, there is a problem in that the optical output characteristics with respect to the drive signal drift timewise due to temperature change and a change with time, resulting in that interference between the codes occurs in the on/off level of the optical output. In order to solve this problem to stably control an operating point of the MZ optical modulator, a control at the operating point is necessary, such that, as shown in FIG. 31, in the case where a curve \u201ca\u201d showing the optical output characteristics with respect to the drive voltage moves to the left to become a curve \u201cb\u201d, a bias voltage by the drive signal is made lower according to the such a change, and in the case where the curve \u201ca\u201d moves to the right to become a curve \u201cc\u201d, the bias voltage is made higher according to such a change. In the conventional modulation method using the NRZ code or the RZ code, a compensation technique has been proposed in Japanese Patent Publication No. 3-251815, for superimposing a low frequency signal onto the drive signal to detect an fluctuation amount and a fluctuation direction of the operating point, and then controlling the bias voltage by feedback to hold the operating point normal.\nUnder the background as described above, also for the new modulation methods such as the optical duo binary modulation method and the CS-RZ modulation method and the like, an operating point control technique for the MZ optical modulator which can correspond to these modulation methods is required.\nAs an operating point control technique for an MZ optical modulator corresponding to the CS-RZ modulation method, there is the one disclosed in Japanese Unexamined Patent Publication No. 2001-119344. In this conventional technique, as shown in FIG. 32, a transmission section of the CS-RZ modulation method comprises: a transmission light source 110 that generates continuous light; an MZ optical modulator 120 on a data modulation side, that is driven in accordance with a drive signal obtained by adjusting, by a drive circuit 121, a data signal corresponding to the NRZ modulation method so that the amplitude thereof becomes V\u03c0, and an MZ optical modulator 130 on a clock modulation side, that is driven in accordance with a drive signal obtained by adjusting, by a drive circuit 131, a clock signal having a frequency of \u00bd the bit rate of the data signal so that the amplitude thereof becomes 2V\u03c0.\nIn such a transmission section of the CS-RZ modulation method, by adjusting a crossing point in an output from the optical modulator 120 on the data modulation side, to an optical extinction phase of the optical modulator 130 on the clock modulation side, an optical signal having the same waveform as for when the RZ code is used is output as an optical output. Therefore, for the optical modulator 120 on the data modulation side, since the drive signal of the NRZ code form with amplitude V\u03c0 is used in the same manner as for the aforementioned case shown in FIG. 31, the operating point can be optimized by applying the conventional operating point control method. On the other hand, for the optical modulator 130 on the clock modulation side, for example as shown in FIG. 33, since the waveform of the drive signal is a sine wave with amplitude 2V\u03c0, it becomes necessary to control the operating point so that the amplitude of the drive signal at the time of optical extinction becomes zero.\nIn addition, for the operating point control method of the MZ optical modulator corresponding to the optical duo binary modulation method, the technique disclosed in Japanese Unexamined Patent Publication No. 2000-162563 is known. In this conventional technique, for example as shown in FIG. 34, different to the case of the modulation method using the NRZ code or the RZ code, a drive signal corresponding to the optical duo binary code, being a differential code with amplitude 2V\u03c0, is supplied to an MZ optical modulator 210. This drive signal is superimposed with an operating point control signal of low frequency generated by an oscillator 212. A part of the optical output from the MZ optical modulator 210 is branched by an optical coupler 220, and then sent to a phase comparator 223 via an optical receiver 221 and an I/V converter 222. In the phase comparator 223, a phase comparison is performed on the superimposed component contained in the optical output, and the operating point control signal, to thereby detect the fluctuation amount and the fluctuation direction of the operating point. Then, the operating point is optimized by feeding back the result of the phase comparison to a control of a bias supply circuit 225. Furthermore, in the conventional technique described above, for example as shown in FIG. 35, a compensation method has also been proposed, for performing a control such that the frequency of the superimposed component contained in the optical output becomes a maximum, by utilizing the phenomenon where the frequency of the superimposed component becomes twice the frequency of the operating point control signal when the bias voltage coincides with the operating point, to hold the operating point normal.\nConsidering the control of the operating point of the MZ optical modulator on the clock modulation side in the CS-RZ modulation method by applying the conventional operating point control technique as described above, even if an operating point control signal with a low frequency is superimposed on the clock signal, since the waveform of the clock signal is a sine wave, the optical output from the optical modulator only becomes a \u201c1\u201d level (apex A of optical emission) temporally, and there is no point where a \u201c1\u201d level is continuous for the optical output, as in the case where an operating point control signal is superimposed on a rectangular wave data signal. Therefore, the superimposed component with a low frequency corresponding to the operating point control signal rarely appears in the optical output from the optical modulator. Consequently, in the conventional method of superimposing the operating point control signal, it is difficult to detect the fluctuation amount and the fluctuation direction of the operating point in the MZ optical modulator on the clock modulation side.\nAlso, in the method which controls the operating point by detecting that the frequency of the superimposed component contained in the output from the optical modulator becomes twice the frequency of the operating point control signal, in the MZ light modulator-on the clock modulation side, since the point where the optical output becomes a maximum is only temporary as described above, the detection signal becomes faint so that it is difficult to determine the frequency of the superimposed component. Moreover, in the conventional method, there is also a problem in that when the amplitude of the drive signal is shifted from 2V\u03c0, then in the optical output, two bias points exist which generate a superimposed component with twice the frequency of the operating point control signal, and in the case where the operating point is controlled based on either of the bias points as a reference, only output waveforms with deviated duty are obtained, so that normal modulation operation cannot be realized.\nFurthermore, in the transmission section of the CS-RZ modulation method as shown in FIG. 32, in order to obtain a desired optical output, an adjustment or a control is required at several points such as the amplitude and the phase between the two drive signals of differential form corresponding to the clock signal, and the phases between the clock signal and the data signal. Therefore, there is a disadvantage in that when providing the additional configurations for the conventional drive control, this invites a higher cost and a larger size for the transmission section."}
-{"text": "An ATM (Automatic Teller Machine) has been widely used in various commercial banks and for postal saving. The ATM operates in the 24-hour self-help manner, which brings people great convenience.\nIn an ATM automatic banknote adding system, all of cash recycling ATMs, cash deposit ATMs and banknote adding machines share a same kind of banknote boxes. The banknote box itself has an information memory for storing information on banknotes in the banknote box. Each banknote box has a physical ID number which is globally unique, so as to be able to distinguish from each other. The automatic banknote adding system may directly transmit a banknote box which is full of banknotes in a cash deposit ATM to a cash outputting ATM, and the cash outputting ATM automatically gets the information on the banknote in the banknote box, so as to automatically add banknotes. The banknote box may be changed at any moment at will between all banknote processing machines in the automatic banknote adding system, and information on each replacing banknote box is uploaded to a server of a bank through each ATM so as to record corresponding data, thus no man-made count is needed, so that the efficiency of processing banknote is greatly improved.\nPresently, if an ATM is used for cash withdrawal, the ATM may send unqualified banknotes to a recovery part. However, information (including the number and face values of the banknotes) on the banknotes sent to the recovery part cannot be correctly determined. Therefore, it is impossible to accurately obtain the information on the residual banknotes in the banknote box.\nFor example, if a certain banknote is in an unqualified status such as an overlapped banknote or an adhesion banknote) when the banknote is output from a cash recycling ATM, it can be determined that the banknote is in the abnormal status by measuring the thickness or width, and thus the banknote is recovered. However, it is very difficult to determine the number of the recovered banknotes (taking the overlapped banknotes as an example, it is impossible to determine whether three or two banknotes are overlapped; and if there is adhesive substance on one banknote, the banknote may also be mistakenly determined as two banknotes). Therefore, it is unable to accurately track the information on the residual banknotes in the banknote box. To accurately obtain the information on banknotes in the banknote box is an important prerequisite for the ATM to automatically add banknotes."}
-{"text": "The present invention relates to a process for the preparation of a low viscosity alkyl toluene or alkyl xylene sulfonate.\nAlkyl toluene and alkyl xylene sulfonates have utility as surfactants for use in enhanced oil recovery. However, these sulfonates are often highly viscous and therefore can be difficult to handle and transport except in highly dilute form. Since the selection of a particular surfactant which can be utilized in petroleum recovery techniques is based in part upon economic considerations, it would be advantageous to have a method for obtaining a low viscosity alkyl toluene or alkyl xylene sulfonate.\nU.S. Pat. No. 3,957,671 describes a low viscosity detergent acid mix containing alkyl benzene sulfonic acid which is prepared by sulfonating alkyl benzene in the presence of benzoic acid."}
-{"text": "1. Field\nThe following description relates to an equal model conservation technique for a pipeline processor.\n2. Description of the Related Art\nAn instruction is processed at several stages in a pipeline processor. For example, an instruction may be processed at a fetch stage, a decoding stage, an execution stage, and a write-back stage. In the pipeline processor, because a plurality of instructions sequentially pass through the individual stages of the pipeline while overlapping each during execution, programs can be processed at high speeds.\nExamples of a pipeline processor include an equal model processor and a less-than-or-equal (LTE) processor. In the pipeline processor, a maximum latency may be set for each instruction. For example, a value of maximum latency for each instruction may depend on a processor type.\n\u201cMaximum latency\u201d indicates a maximum time period for which an operation corresponding to each instruction is executed and the result of the execution is written in a target register. In the equal model processor, the result of the execution of an instruction is written in the target register after the set maximum latency elapses. Meanwhile, in the LTE model processor, the result of execution of an instruction may be written in the target register before set maximum latency elapses. Accordingly, the equal model processor allows flexible scheduling through multiple assignment of a target register because the result of the execution of an instruction is not written in the target register until maximum latency elapses. Alternatively, in the LTE model processor, because the result of execution of an instruction can be written in a target register before maximum latency elapses, multiple assignment of the target register is limited."}
-{"text": "In the medical field, it is sometimes the case that when a region of interest exists on an image obtained by a given modality, a corresponding region is identified on an image obtained by another modality, and diagnosis is then performed based on the comparison between them. If imaging posture differs between modalities, the shape of the object differs at the times of imaging. This makes it difficult to identify the object. Under the circumstance, attempts have been made to estimate the deformation of an object between modalities (that is, alignment between images accompanying deformation). This makes it possible to estimate the position of a corresponding region based on the position information of a region of interest and to generate an image by deforming one image so as to have the same shape as that of the other image.\nNPL 1 discloses a technique of facilitating the comparison between the shapes of an object before and after its deformation by normalizing the shape of the object accompanying deformation. More specifically, there is disclosed a method of calculating the geodesic distance matrix of the surface shape of an object and normalizing the surface shape by using a multidimensional scaling method using the calculated matrix. This method can normalize the shape of an object which deforms without changing the geodesic distance matrix of the object surface so as to allow direct comparison between the shapes of the object before and after the deformation. This makes it possible to easily compare the deformed shapes of an object accompanying deformation and recognize the object based on its shape."}
-{"text": "The present invention relates to boiling heat transfer surfaces which are provided on a heat exchange wall surface to enhance nucleate boiling and thereby improve the heat transfer action of a boiling liquid such as a halocarbon refrigerant. It is well known that nucleate boiling in a heat transfer apparatus is improved when the heated wall portion through which heat is transmitted from a warm fluid to a boiling liquid has its surface made porous by the formation of cavities in the heated wall portion.\nU.S. Pat. No. 3,384,154 teaches that the radii of the cavities provided in the heated wall portion can be controlled by the size of the sintered powder coatings utilized to make the nucleate boiling surface.\nU.S. Pat. No. 3,906,604 controls the nucleate boiling cavity size by machining the metal substrate to create protrusions which are frictionally rubbed and stretched to form the cavities.\nU.S. Pat. Nos. 3,696,861 and 3,768,290 utilize bent-over tube fins to create elongated nucleate boiling cavities.\nGerman Publication No. 2,510,580 discloses heat pipe surfaces wherein the surface is first mechanically roughened by wire brushing to form close uniform scratches and then electroplated to form nucleate boiling cavities.\nU.S. Pat. No. 4,018,264 discloses a boiling heat transfer surface wherein the nucleate boiling cavities are made by electrodepositing copper dendrites on a copper surface at a high current, copper plating the dendrites or nodules and then rolling the dendritic surface to partially compact the nodules.\nThe present invention relates to a heat transfer surface having improved cavities to promote nucleate boiling; an improved method for preparing such cavities; and an improved method of nucleate boiling which utilizes the improved cavity structures. The present method for preparing the improved nucleate boiling cavities comprises deforming the metal substrate to form indentations in the metal in the form of non-continuous pits or craters and then electrodepositing a metal on the periphery of the pits or craters at a high current density and then at a normal plating current, from a special bath, to provide honeycomb-like surfaces which optimize nucleate boiling augmentation for particular fluids, especially halocarbon refrigerants utilized in refrigerator systems.\nThe cavity diameters in the plated coatings are directly related to the size of the pits or craters formed on the metal substrate. In the case of sandblasting, the metal deformation is made by grits having a mesh size of 10 to 200. In the case of deforming by knurling, the cavity diameter is determined by the size and number per square inch of the projections on the knurling tool. A pre-plated surface having from 100 to 350 (10,000 to 122,500 per square inch) knurls per inch is preferred in making the present heat exchange surface. It has been found that a knurling tool having 240 knurls per square inch provides pits or craters which, when electroplated, form excellent nucleate boiling cavities in the heated wall surface of a heat exchanger using halocarbon refrigerants such as R-12.\nWhen the heat exchange surface is deformed by sandblasting, silicon carbide of 10-200 mesh, glass beads of 10-200 mesh and ground slag of 60 mesh have been used; silicon carbide of 60 mesh size is preferred for use in R-12."}
-{"text": "1. Field of the Invention\nThe present invention relates to a multi-layered circuit substrate for a semiconductor device. Further, the present invention relates to a method of manufacturing the multilayered circuit substrate. Furthermore, the present invention relates to a semiconductor device using such a substrate and a method of producing the same. More particularly, the present invention relates to a multilayered circuit substrate for a semiconductor device having a semiconductor element mounting face on which a plurality of semiconductor elements can be mounted being arranged in the plane direction on one side of the multilayered circuit substrate composed by laminating a plurality of conductor patterns. Further, the present invention relates to a method of producing the multilayered circuit substrate, a semiconductor device, and a method of producing the same.\n2. Description of the Related Art\nRecently, there has been provided a semiconductor device referred to as a system-in-package (SIP) in which a plurality of semiconductor elements are mounted on one circuit substrate. This semiconductor device will be referred to as \u201cSIP\u201d, hereinafter.\nA size of the aforementioned SIP is appropriate to be handled as a chip. An example of SIP is shown in FIGS. 7(a) and 7(b). SIP 100 shown in FIG. 7(a) is composed in such a manner that semiconductor elements 104a, 104b, 104c, the functions of which are different from each other, are mounted in the plane direction on a semiconductor element mounting face formed on one face of a piece of multilayered circuit substrate 102.\nAs shown in FIG. 7(b), the multilayered circuit substrate 102 forming the above SIP 100 is composed of resin layers 102a, 102b, 102c, 102d which are made of insulating resin, and the conductor patterns 106, 106, . . . and the via holes 108, 108, . . . are formed being laminated on these resin layers 102a, 102b, 102c, 102d. \nOn one face of the multilayered circuit substrate 102, there are provided connection pads 110, 110, . . . , from which the connecting faces to be connected with the electrode terminals of the semiconductor elements 104a, 104b, 104c are exposed. On the other face of the multilayered circuit substrate 102, there are provided external connection pads 114, 114, . . . , from which the attaching faces on which the solder balls 112, 112, . . . are attached are exposed.\nThe connection pads 110, 110, . . . and the external connection pads 114, 114, . . . are electrically connected with each other by the conductor patterns 106, 106, . . . and the via holes 108, 108, . . . which are formed and laminated on the resin layers 102a, 102b, 102c, 102d. \nOne face and the other face of the above multilayered circuit substrate are covered with the protective films 116, 118 made of solder resist except for the connection pads 110, 110, . . . and the external connection pads 114, 114, . . . .\nIn this connection, the potting resin 120 is filled between the semiconductor element mounting face, which is formed on one face of the multilayered circuit substrate 102, and the semiconductor elements 104a, 104b, 104c mounted on the semiconductor element mounting face.\nWhen SIP 100 shown in FIGS. 7(a) and 7(b) is used, the size of SIP can be reduced to be smaller than the size of a device in which a plurality of semiconductor devices, each semiconductor device having a single semiconductor element, are used. Further, it is possible to reduce a conductor distance between the semiconductor elements 104a, 104b, 104c. Therefore, a signal can be sent and received between the semiconductor elements at high transmission speed.\nHowever, SIP 100 shown in FIGS. 7(a) and 7(b) is composed of a multilayered circuit substrate 102 actually made of resin. Therefore, rigidity of the multilayered circuit substrate 102 is not sufficiently high.\nTherefore, the present inventors made investigations into SIP 200 shown in FIGS. 8(a) and 8(b). As shown in FIG. 8(b), the multilayered circuit substrate 202 composing this SIP 200 is formed by laminating the resin layers 202a, 202b made of insulating resin on which conductor patterns 106, 106 . . . and the via holes 108, 108 . . . are formed. On one face of the multilayered circuit substrate 200, there are provided connection pads 110, 110, . . . , from which the connecting faces to be connected with the electrode terminals of the semiconductor elements 104a, 104b, 104c are exposed.\nOn the other face of the multilayered circuit substrate 202, the metallic plate 204, which is a plate-shaped member having a rigidity higher than that of the multilayered circuit substrate 202, is bonded by the adhesive layer 206 made of insulating resin.\nSince metallic plate 204 is joined to the other face of the multilayered circuit substrate 202, SIP 200 and other electronic parts are electrically connected with each other by a lead frame as shown in FIGS. 8(a) and 8(b). Specifically, an end portion of each inner lead 300, 300, . . . of the lead frame is connected with an exposure face of each external connection pad 208, 208, . . . formed along the outer edge of the multilayered circuit substrate 202 by the wire 302.\nIn this SIP 200, the highly rigid metallic plate 204 is joined to the other face of the multilayered substrate 202. Therefore, rigidity of this SIP 200 is actually enhanced and higher than that of the multilayered circuit substrate 102 made of resin shown in FIGS. 7(a) and 7(b).\nHowever, since the coefficient of thermal expansion of the multilayered circuit substrate 202 actually made of resin and that of the metallic plate 204 are different from each other, stress is generated between them. Cracks tend to be caused on SIP 200 by the thus generated stress.\nSince the multilayered circuit substrate 202 shown in FIGS. 8(a) and 8(b) is successively laminated from the resin layer 202a provided on the metallic plate 204 side. Therefore, on a surface of the resin layer 202b laminated on the resin layer 202a, especially on a surface of the resin layer 202b corresponding to the via hole 108 formed on the resin layer 202a, a concave surface, or indentation, tends to be formed as shown in FIG. 9. Therefore, the surface of the resin layer 202b tends to become irregular.\nWhen the connection pad 110 to be connected with the electrode terminal of the semiconductor element is formed on the irregular face of the resin layer 202b as shown in FIG. 9, an exposed face of the connection pad 110 is formed into an irregular face, following the irregular face of the resin layer 202b. \nIn the case where the exposed face of the connection pad 110 is formed into an irregular face, when a semiconductor element is mounted on a semiconductor element mounting face of the multilayer circuit substrate 202, an electrode terminal of the semiconductor element does not come into contact with the exposed face of the connection pad 110, which causes an imperfect contact and reliability of the finally obtained SIP 200 is deteriorated.\nAlthough the irregular face formed on the resin layer 202b of the multilayer circuit substrate 202 can be flattened by means of polishing, it is necessary to add a polishing process to the conventional manufacturing process of SIP 200, which raises the manufacturing cost of SIP 200. For the above reasons, it is preferable that the semiconductor element mounting face of the multilayered circuit substrate 202 is flattened without adding the polishing process."}
-{"text": "It is well known within the electronics industry that lead coplanarity is essential to the successful placement of fine-pitch surface mount devices (SMD's). Non-coplanar leads, i.e. those bent up or down with respect to the body of the component, typically result in opens or poor solder joints on the printed circuit board (PCB). It is desirable to reject any components that do not meet coplanarity requirements prior to placement on the printed circuit board. The present industry standards require that all lead tips be coplanar within a 0.004 inch band, forcing inspection equipment to measure coplanarity with .+-.0.002 inch accuracy. Although some non-contact 3-D inspection techniques are capable of providing this degree of accuracy, a number of system constraints prohibit their implementation for this application. These constraints include but are not limited to:\n(1) large part sizes,\n(2) establishment of a reference (datum) plane,\n(3) mechanical space and packaging limitations within the pick and place machine, and\n(4) real-time inspection to maintain high pick and place throughput rates.\nOther techniques, particularly \"implied\" techniques based around passive triangulation (i.e. estimation of lead height via measured lead length), have severely limited depth resolution. They may also have other drawbacks such as confusion between lead length error and vertical lead displacement as well as ambiguous polarity information (the inability to differentiate between leads bent up or down with respect to the plane of the part).\nIt is accordingly an object of the present invention to overcome these constraints while providing the necessary degree of accuracy for proper inspection of lead coplanarity. The optical setup inherent in the disclosed system provides a true, direct measurement of lead height. It produces sufficient depth resolution to measure lead coplanarity to an accuracy of .+-.0.002 inch using conventional sensors and image processing hardware.\nAn advantage of the disclosed system is its ability to handle a range of SMD components including PQFP, SOIC, PLCC, SOJ, and TapePak with pitch values down to 15 mils. It will accommodate 25-mil pitch parts up to 1.25 inches square in a single field of view. Parts with a larger number of leads or finer leads can also be inspected by breaking up the image into multiple fields of view, typically two or four. If necessary, large fine-pitch devices may be inspected through the use of high resolution image sensors such as line scan or non-RS-170 area cameras, though at present this requires additional complexity for image acquisition.\nOther objects and advantages of the present invention will become apparent from the following portion of the specification and from the following drawings which illustrate in accordance with the mandate of the patent statutes a presently preferred embodiment incorporating the principles of the invention."}
-{"text": "1. Field of the Invention\nThe present invention relates to an image signal compressing device adapted for use in a facsimile machine and an electronic filing machine.\n2. Description of the Prior Art\nRecently, a device has been developed which, when used in a facsimile machine or an electronic filing machine, operates to convert an image of intermediate gradation into a binary image signal according to a dither process.\nIt has, however, been found that, when the binary image signal, converted by the above described device, is compressed according to an MR code or a modified Hoffman (MH), which is a coding (compressing) system standardized in the facsimile machines, the compression ratio tends to be considerably lowered."}
-{"text": "The present invention is directed toward optical communication system. In particular, optical communication systems providing dispersion slope compensation.\nOptical signals transmitted in a fiber optic communication system typically constitute a series of pulses of digital information. Although the pulses are usually at a single nominal wavelength, each pulse is actually composed of different spectral components. These spectral components propagate through the transmission fiber at different speeds with higher frequency components traveling slower than lower frequency components. This effect, known as xe2x80x9cchromatic dispersionxe2x80x9d, can result in spectral components of one pulse arriving at a receiver at substantially the same time as a succeeding pulse, thereby causing degraded receiver sensitivity. Chromatic dispersion becomes increasingly pronounced at higher bit rates, e.g. 10 Gbit/sec. such as those associated with synchronous optical network (SONET) OC192 transmission speeds. Unique fibers have been developed to offset or compensate chromatic dispersion. These fibers, referred to as dispersion compensated fiber or DCF, are commercially available from Corning Inc. and Lucent Technologies, for example. DCF fibers are typically specified with a dispersion coefficient having an opposite sign to that of the transmission fiber. The net dispersion experienced by the transmitted optical signal is thus the sum of the dispersion accumulated through the transmission fiber and the DCF. If the product of dispersion coefficient times length in the DCF is equal in magnitude and opposite in sign to that of the transmission fiber, the net dispersion is zero, and the dispersion is said to be compensated.\nDCF typically provides compensation at a particular wavelength. As a result, single wavelength optical transmission systems incorporating DCF and operating at the specified wavelength can have relatively low error rates, even at OC192 transmission speeds.\nIn order to further increase fiber capacity, wavelength division multiplexed (WDM) systems have been developed for carrying multiple wavelengths on a single fiber. At OC192 rates, each wavelength or channel of a WDM system must be dispersion compensated. Typically, the dispersion coefficient of the transmission fiber, as well as that associated with DCF, varies with wavelength; the relationship of dispersion and wavelength being referred to xe2x80x9cdispersion slopexe2x80x9d. Since the dispersion slope associated with the transmission fiber and DCF may not be identical, not all channels may be adequately dispersion compensated.\nTheoretically, one approach to providing spectrally uniform dispersion compensation involves fabricating DCF uniquely tailored to have a slope of the same magnitude, but opposite sign, as the transmission fiber. As a practical matter, however, such DCF would be difficult to fabricate with precisely the required slope and magnitude.\nAlternatively, since the amount of dispersion compensation depends in part on the length of the DCF, varying lengths of DCF can be provided in a WDM system downstream from the point where individual wavelengths have been separated. Each channel can then be compensated individually prior to detection at a photodetector. An additional segment of DCF, however, must be provided for each channel, thereby complicating system design and increasing costs.\nConsistent with the present invention, an optical device is provided comprising an optical splitter having an input port and a plurality of output ports, said input port being configured to receive a plurality of optical signals, each at a respective one of a plurality of wavelengths. The optical signals are output from each of the plurality of output ports, but in attenuated form. The optical device further includes a plurality of in-fiber Bragg gratings, each of which being coupled to a respective one of the plurality of output ports. The in-fiber Bragg gratings are configured to reflect a selected one of the plurality of optical signals and compensate for a dispersion associated with the selected one of the plurality of optical signals."}
-{"text": "This disclosure relates to the polishing of magnetic, optical, semiconductor and silicon wafers and more particularly, to polishing compositions and methods for polishing and planarizing silicon wafers.\nThe semiconductor industry uses interconnect metals in forming integrated circuits on semiconductor wafers. These interconnect metals are preferably non-ferrous metals. Suitable examples of such non-ferrous interconnects are aluminum, copper, gold, nickel, and platinum group metals, silver, tungsten and alloys comprising at least one of the foregoing metals. These interconnect metals have a low electrical resistivity. Copper metal interconnects provide excellent conductivity at a low cost. Because copper diffuses into many dielectric materials, such as silicon dioxide or doped versions of silicon dioxide, integrated circuit fabricators typically apply a diffusion barrier layer to prevent the copper diffusion into the dielectric layer. For example, barrier layers for protecting dielectrics include, tantalum, tantalum nitride, silicon nitride, tantalum-silicon nitrides, titanium, titanium nitrides, titanium-silicon nitrides, titanium-titanium nitrides, titanium-tungsten, tungsten, tungsten nitrides and tungsten-silicon nitrides.\nIn the manufacturing of semiconductor wafers, polishing compositions are used to polish semiconductor substrates after the deposition of the metal interconnect layers. Typically, the polishing process uses a \u201cfirst-step\u201d slurry specifically designed to rapidly remove the excessive interconnect metal. The polishing process then includes a \u201csecond-step\u201d slurry to remove the barrier layer. The second-step slurry selectively removes the barrier layer without adversely impacting the physical structure or electrical properties of the interconnect structure. In addition to this, the second step slurry should also display low erosion for dielectrics.\nWhile the \u201ctwo-step\u201d slurry generally performs satisfactorily in present conditions, it is desirable to have polishing compositions that can be used for planarizing the metal interconnect without substantial dishing of the interconnect metal and scratching of the wafer surface.\nPreparation and evaluation of chemically modified abrasives for chemical mechanical polishing by Partch et al. in an abstract presented at Materials Research Society (MRS) meeting in 2001, discloses coating individual particle cores with a layer of organic or inorganic material. However, Partch et al., disclosed manufacturing the polishing composition with abrasive particles having a large particle size distribution that resulted in scratching and other forms of surface damage.\nU.S. Pat. Pub. No. 2004/0060502 to Singh discloses a plurality of particles that may be coated by weakly or strongly adsorbing surfactants or polymer additives to a core particle. The polymers thus adsorbed onto the abrasive particles can desorb from the particles during the polishing process to alter polishing characteristics and possibly result in damage to the polished surface.\nThere thus remains an unsatisfied demand for polymer-coated particles suitable for chemical mechanical polishing of magnetic, optical, semiconductor or silicon substrates."}
-{"text": "There is a great variety of tools for picking up and manipulating small objects, arranging for example from the simplest tongs and magnetized screwdrivers to extremely complex remote control arms. Many of these tools, such as tongs or pliers, require that the operator maintain his grip on the work piece being handled. Others, such as certain forceps and locking pliers, can be locked in a closed position on the piece being handled. Many of the locking manipulative tools are relatively large, clumsy, otherwise not particularly well-adapted for manipulation of small electrical components, or too expensive for the task. As background, see for example U.S. Pat. No. 4,605,256 which is directed to a locking tool for manipulation of electronic components, and U.S. Pat. No. 3,896,533, which is directed to a tool for inserting and removing circuit components.\nElectrical outlets are frequently crowded environments in which it can be dangerous or difficult to handle components with the fingers. Accordingly, it is desirable to have tools which are small or portable, easily operable and manipulated by one hand, and which can grasp electrical components firmly, capturing them so that they can be manipulated in any orientation for installing and testing or removing electrical outlets from junction boxes. The tool of the present invention admirably fulfills these requirements. It is simple, easily operable by one hand, can be interconnected with the work piece being handled, and is able to capture and retain the electrical work piece being handled until the operator chooses to disconnect the tool from the work piece."}
-{"text": "Electronic compasses are not new in the art. Both flux gate and Hall effect circuitry have been used to create suitable directional displays for a variety of navigational tasks.\nFor example, U.S. Pat. No. 4,402,142, issued to Robert C. Dinsmore on Sep. 6, 1983, describes an electronic compass for automobiles utilizing a series of Hall effect sensors. However, the circuits used in that invention have proven to require significant amounts of current, rendering the devices relatively impractical for portable use. Because of the current demands of the Hall effect sensors, amplifiers, Schmitt triggers and display drivers, present Hall effect devices typically consume current at the rate of approximately 200 milliamperes per hour. Extremely large batteries would be necessary to power these devices for extensive periods of time.\nThe within invention eliminates substantial current consuming portions of the typical Hall effect or flux gate-type compass, adds a \"chopper\" circuit which effectively shuts the compass electronics off during most of its operating cycle, and a vibration sensing device which shuts off electrical power to the compass when it is motionless for a specified time."}
-{"text": "1. Field of the Invention\nThis invention relates to an apparatus for starting and operating a discharge lamp, using an inverter circuit.\n2. Description of Background\nAn apparatus for starting and operating a discharge lamp, equipped with an inverter circuit has generally been used so as to achieve a small-sized unit with an improved efficiency. Japanese Patent Disclosure (KOKAI) No. 57-141895 discloses a conventional apparatus for starting and operating a discharge lamp. According to the invention, an AC power source is connected to a rectifier circuit including a rectifier and a smoothing capacitor, the output of the rectifier circuit being connected to an inverter circuit. The inverter circuit includes a parallel resonance circuit of an oscillation transformer and capacitor connected in parallel with a primary winding of the oscillation transformer. The resonance circuit is driven by a switching transistor. Between the base and the emitter of the transistor, a series circuit of a coupling capacitor and feedback coil magnetically coupled to the oscillation transformer is connected as a feedback circuit, via a diode. A series circuit of a diode and resistor is connected in parallel with the series circuit of the coupling capacitor and feedback coil, and a trigger resistor is connected in series with both the series circuits.\nA discharge lamp is connected to a secondary coil of the oscillation transformer, via a current-limiting coil, and a starting circuit is comprised of a series circuit of diodes and silicon symmetrical switch element, and is connected in parallel with a discharge lamp.\nAccording to the aforementioned apparatus for starting and operating a discharge lamp, since the starting circuit includes a series circuit of diodes and a silicon symmetrical switch element, a high-voltage pulse to be generated is regulated by the diode characteristic. It is, therefore, not possible to obtain an adequate high-voltage pulse."}
-{"text": "This invention relates in general to chucks for holding tools and in particular to a new and useful power operated chuck which includes radially movable clamping jaws.\nSuch power chucks, also termed wedge hook chucks, are known in a variety of designs and have proved satisfactory in practice. To equalize the centrifugal forces produced in the clamping jaws at high speeds and reduce an inwardly directed static clamping force or considerably augment an outwardly directed clamping force, it is known to provide balancing elements which are guided in the chuck body as counterweights and operatively connected to the jaws through hinged levers. The principal drawback of these constructions is that the static clamping forces acting on the clamped workpiece before and after a machining operation are not equal to each other and differ sometimes considerably. That is, as the chuck rotates at high speeds, the clamping force transmitted from the clamping piston through the wedge hooks to the jaws becomes almost counterweighted by the produced centrifugal forces and the elasticity of the coupled parts. Then, a part of the clamping force is produced by the balancing elements which are connected to the jaws through reversing levers. The clamping piston coupled to the jaws through the mating surfaces of the wedge hooks is thereby relieved and displaced, even though very little, by the action of a connected servo-mechanism which continues to be effective. In addition, the centrifugal forces expand the chuck body, so that the bore receiving the clamping piston is enlarged, with the result that the no longer backed-up wedge hooks can yield, which causes a further displacement of the clamping piston. However, after a machining operation is terminated, the chuck body contracts again, whereby the wedge hooks are pressed inwardly and the clamping force on the workpiece is augmented since the clamping piston does not retract, it is held fast in its displaced position. In consequence, the clamping force may increase considerably and often leads to a deformation of the workpiece."}
-{"text": "1. Field of the Invention\nThe present invention relates to integrated circuit lead frames and methods of production thereof. In particular, this invention relates to plastic lead frames with a conductive coating or material contained therein used for packaging integrated circuits and methods of manufacturing the same.\n2. State of the Art\nIntegrated circuit (IC) chips are enclosed in plastic packages that provide protection from hostile environments and enable electrical interconnection to printed-circuit boards. During a manufacturing process, the IC chip is typically attached to a die paddle of a conventional lead frame or suspended from the lead fingers of a leads-over-chip (LOC) lead frame using an adhesive such as epoxy or double-sided tape, and subsequently encapsulated with a dense and rigid plastic by a transfer molding process. In essence, the lead frame forms the backbone of the molded plastic IC package.\nLead frames typically perform many functions such as: (1) a holding fixture that indexes with tool-transfer mechanisms as the package proceeds through various assembly operations, (2) a dam that prevents plastic from rushing out between leads during the molding operation, (3) a chip attach substrate, (4) a support matrix for the plastic, and (5) an electrical and thermal conductor from chip to board.\nTraditionally, lead frames are fabricated from a strip of sheet metal by stamping or chemical milling operations. There are many different metal alloy compositions which are commercially available for producing lead frames. For example, Rao R. Tummala and Eugene J. Rymaszewski, \"Microelectronics Packaging Handbook,\" Table 8-4, 1989, provide 16 different alloys available from 9 different manufacturers. Lead frame material selection depends on many factors such as cost, ease of fabrication, strength, thermal conductivity, and matched coefficient of thermal expansion (CTE). A close match of CTE between the silicon die and the lead frame is required to avoid chip fracture from differential expansion rates.\nThe most widely used metal for lead frame fabrication is Alloy 42 (42% Nickel-58% Iron). Alloy 42 has a CTE near silicon and good tensile strength properties. The disadvantage of Alloy 42 is that it has low thermal conductivity. Since the lead frame is the main conduit by which heat flows from the chip to the environment and printed circuit board, this can have a profound effect on the package thermal resistance after prolonged device operation.\nA layered composite strip, such as copper-clad stainless steel, was developed to emulate the mechanical properties of Alloy 42 while increasing thermal conductivity. However, copper-clad stainless steel is somewhat more expensive to manufacture than Alloy 42. When manufacturing copper-clad stainless steel lead frames, the cladding is accomplished by high-pressure rolling of copper foil onto a stainless steel strip, followed by annealing the composite to form a solid-solution weld. While copper alloys provide good thermal conductivity and have a CTE near that of low-stress molding compounds, there is a substantial CTE mismatch with respect to silicon.\nWhile numerous alloys have been developed to solve problems with thermal conductivity, CTE mismatch, and strength, other important factors such as ease of fabrication and cost have not improved as readily.\nConventional methods for making lead frames for integrated circuit devices are described in U.S. Pat. No. 3,440,027. The use of a plastic support structure in a method of forming metal lead frames is described in U.S. Pat. No. 4,089,733 (hereinafter the \"'733\" patent). The plastic support structure of the '733 patent solves the problem of deformed and misaligned lead fingers resulting from stress during the bonding process by supporting the lead fingers with a plastic structure. However, the '733 patent requires a metal lead frame in addition to the plastic support structure with its attendant costs. A method of manufacturing multi-layer metal lead frames is disclosed in U.S. Pat. No. 5,231,756 (hereinafter the \"'756\" patent). The '756 patent provides an improvement in aligning power and ground planes for use in a multi-layer lead frame where such planes are necessary. However, the number of steps required to manufacture such multi-layer lead frames will not solve the problem of decreasing costs. In short, none of these related art appear to disclose methods of producing low-cost lead frames made from materials not structurally based on metal.\nSince packaged ICs are produced in high volumes, a small decrease in the cost per packaged IC can result in substantial savings overall. Accordingly, there is a need in the industry for a low-cost plastic lead frame with suitable characteristics for IC packaging."}
-{"text": "In a storage device using a semiconductor memory or the like, error correction processing is performed in data decoding processing. There is a chien search as a method of specifying error locations of data in the error correction processing. In the chien search, exclusive-OR operations multiplying \u03b1 to the power, which multiply coefficients of an error location polynomial (t-order or less) except for a constant term by \u03b1, \u03b12, . . . , \u03b1t, and exclusive-OR operations, which obtain the sum of the respective results of the operations on a Galois field, are performed to specify error locations by searching roots. Since error locations are determined in every cycle in a chien search circuit, all the exclusive-OR operations are performed in one cycle. Meanwhile, it is possible to improve error correction capability by increasing the order of the Galois field or the order of the error location polynomial in the chien search circuit, but the number of stages of the exclusive-OR operations is increased."}
-{"text": "Conventionally, there is a medicine dispensing device for dispensing blister packages having a configuration in which a blister package is conveyed with the help of a grip unit and cut by a cutter mechanism to enable retrieval of the required quantity of the package sheet (see Patent document 1, for example).\nAs another medicine dispensing device for dispensing blister packages, there is one having a configuration that enables retrieval of a blister package loaded inside a medicine cassette by sucking it with a suction member (see Patent document 2, for example).\nHowever, in the medicine dispensing device disclosed in Patent Document 1, a grip unit is used for retrieving not only a fraction of the blister package but also one full blister package, and therefore, there is a problem of poor retrieving efficiency. Further, since the blister packages are stacked vertically, the space occupied in the height direction increases, and there is a problem that the quantity and type that can be accommodated is limited.\nIn the medicine dispensing device disclosed in Patent Document 2, it is not possible to dispense fraction of a blister package."}
-{"text": "This invention relates to pneumatic tires and more particularly to a pneumatic tire having a wide reinforcing breaker or belt extending into an annular protuberance in the buttress region of the tire.\nThe tread reinforcement component commonly referred to as a belt or breaker is disposed between the tread and the carcass in the crown region of the tire. Typically, a breaker or belt is constructed of one or more plies of rubberized cord fabric. Each cord in the fabric is generally of a substantially non-extensible material such as metallic wires or strands, glass filaments, or rayon filaments. The cords within the breaker plies are oriented parallel to each other and substantially parallel to the plane of the beads and thus circumferentially of the tire, or at a small bias angle. In a single ply breaker construction, the small bias angle can be 0.degree.. Where the breaker is a multi-ply structure, however, the small bias angle is in the range of about 15.degree.-30.degree. with respect to the median equatorial plane of the tire, and may be oppositely disposed in successive plies.\nThe development of a reinforcing belt or breaker for use in pneumatic tires added a number of benefits to tires including considerable improvement to tire durability. On the other hand, the breaker also introduced some new problems of its own. The lateral edges of a breaker in a typical tire, due to their location between the tire tread and carcass at a position of continually changing stress, tend to separate from the carcass or themselves. Improper positioning of the breaker results in the uniformity problems of increased lateral force variation and conicity.\nAnother benefit of a reinforcing breaker is improved tread wear due to the greater stiffness imparted to the tire tread. Of course, this benefit is enhanced by utilizing a \"flat\" tread having no convex cross-sectional curvature so that the relative movement between the center portion of the tread and the axially outward portions is minimized thereby resulting in an even wear pattern. A serious drawback of a tread designed to have no convex cross-sectional curvature lies in the fact that a tire which is vulcanized with such a shape, when cured, may inflate to a concave tread cross-section rather than a flat cross-section. Elimination of this problem can be accomplished by constructing and vulcanizing the tire such that the tread has at least a small convex cross-sectional curvature, that is, a large crown radius ratio, when the cured tire is inflated to its rated inflation pressure.\nIn the detailed description of the invention following hereinafter, the cross-sectional curvature of the tread and breaker of the tire are discussed in terms of \"crown radius ratio\" and \"crown radius.\" The term \"crown radius\" as used herein either with respect to the tread or breaker, means the radius of the arc of a circle which best approximates the arc of curvature of the tread surface or the breaker as the case may be of a tire. The crown radius is commonly determined by the formula ##EQU1## where \"leg set\" is the axial distance measured along a line tangent to the tread or breaker at its midpoint between the shoulders of the tread or edge of the breaker, or predetermined points axially equidistance inward of the shoulders or breaker edges. Shoulder drop is the perpendicular distance from the aforementioned tangent line to the tread shoulders or said predetermined points. In the case of a breaker, shoulder drop may be considered as the perpendicular distance between the line tangent to the breaker at its midpoint and the edges of the breaker. The \"crown radius ratio\" is the ratio of the crown radius of either the tread or breaker to the nominal cross-sectional width of the tire and as used herein is expressed in terms of percent.\nEither a belted radial or a belted bias tire construction is suitable for use in the present invention. The radial type tire is, however, the preferred embodiment. A radial tire is inclusive of various tire constructions which typically comprise a carcass or body having one or more reinforcement plies of cord fabric extending from bead to bead wherein the cords in each ply are substantially radial in orientation, that is, the cords are oriented substantially normal to the beads and the crown centerline of the tire. In a single ply radial tire construction, the carcass cords normally have a 90.degree. bias angle, that is, in the unshaped carcass they extend perpendicular to the planes of the beads. In a two-ply radial tire construction, the cords in each carcass ply are parallel to the cords in the other carcass ply. However, the cords in each ply may be oriented at oppositely disposed angles of 70.degree. or more with respect to the median equatorial plane of the tire and thus the angle between cords in different plies is between 0.degree. and 40.degree.. In more than two ply radial ply tire constructions, similar cord arrangements in successive carcass plies is usually employed.\nWith regard to the prior art, U.S. Pat. No. 3,450,182 discloses a tire having a flat tread cross-section, wide protuberances extending beyond the width of the tire sidewalls, a tread surface in normal engagement with the ground between the lateral edges of the protuberances and a breaker extending laterally beneath the tread to positions adjacent the edges of the protuberances. This patent does not, however, disclose a breaker having its lateral edges removed from the stress area of contact by the tread with the ground or a convex cross-sectional tread curvature facilitating the manufacture of a tire. Moreover, the extreme lateral extension of the protuberances beyond the tire sidewall will result in a high degree of tread shoulder wear and breaker separation due to centrifugal force moving the protuberances radially outward during tire rotation.\nIn U.S. Pat. No. 2,477,754, an aircraft tire is shown in which a tread and breaker have a substantial convexly curved cross-section when the tire is inflated and the breaker has edges positioned within protuberances extending from the side of the tire. When this tire is loaded (but not impacted), the tread-ground contact area is in the crown area of the tread only and consequently the tread wear and traction of the tire will be relatively poor. Moreover, when the tire is subject to the normal type of impaction and very heavy loading which aircraft tires receive, the lateral edges of the breaker rotate through the tread-ground contact area to cause edge separation of the breaker. A further defect in the tire of U.S. Pat. No. 2,477,754 is that the circumferential breaker cords along the lateral sides of the breaker will carry the entire breaker load under deflection so that these breaker cords will break after a small amount of such use."}
-{"text": "Collapsible tube squeezing devices for dispensing the contents of tubes containing toothpaste, hair creme, hand lotion, shoe polishing creme, glue, and the like are well known in the art. A review of the prior art devices shows that while such devices exist they have not received substantial user acceptance due to their expense, size, lack of universality, difficulties with operation and the necessity to use both of the operator's hands in emptying the contents of collapsible tubes.\nThe advent of collapsible tubes formed of plastic, rather than the well known metallic type collapsible tubes, has also presented additional problems which have rendered many pre-plastic tube squeezing devices unsuitable for their intended purposes particularly in regard to removal of the contents of the crown portion of such tubes. Crown portions of plastic collapsible tubes have a tendency to be substantially more rigid than metallic tubes which factor coupled with the resiliency of plastics often employed in the manufacture of plastic collapsible tubes necessitates that special measures be undertaken in order to completely empty such tubes."}
-{"text": "1. Field of the Invention\nThe present invention relates to an information processing apparatus, a control method therefor, and a computer-readable storage medium storing a program for implementing the control method.\n2. Description of the Related Art\nIn recent years, as image forming apparatuses have become increasingly sophisticated, it has become common to store user data in a nonvolatile device such as an HDD (hard disk drive) or flash memory and read and use the user data when using a certain function.\nIn this case, when a nonvolatile device in which user data is stored fails, it becomes impossible to read the user data.\nAs a countermeasure against data loss caused by such a nonvolatile device failure, there is a mirroring function (RAID1) of constantly backing up data using two HDDs and providing notification when the HDDs fail.\nThere is also a failure predicting function of, before an HDD fails, using S.M.A.R.T. (Self-Monitoring, Analysis and Reporting Technology) information on the HDD to provide notification of a failure prediction indicating that a failure of the HDD is anticipated, thus urging a user to back up data.\nIn a case where, out of the mirroring function and the failure predicting function described above, only the mirroring function is enabled, data is constantly backed up, and hence a service engineer has only to replace the HDD after the HDD fails.\nUsing an HDD to the limit as described above has a merit of, for a user, eliminating downtime of an image forming apparatus caused by occurrence of an error, and a merit of, for the service engineer, reducing maintenance costs.\nOn the other hand, in a case where only the failure predicting function is enabled, upon receiving notification of a failure prediction, a service engineer needs to replace an HDD after backing up user data so as to ensure the user data. This failure prediction, however, is merely a prediction, and hence the HDD may fail immediately after notification or may be usable for several years after notification.\nIn a case where both the mirroring function and the failure predicting function are enabled, when notification of a failure prediction using the failure predicting function is provided, notification of an error is provided to a user although an HDD can still be used.\nAs a result, it is necessary to replace the HDD of an image forming apparatus for which notification of the error has been provided, resulting in undesired downtime and cost of maintenance by a service engineer.\nMoreover, when both the mirroring function and the failure predicting function are enabled, this means that an image forming apparatus is equipped with two HDDs. In this case, when notification of a failure prediction with respect to both the two HDDs is provided, it is necessary to notify a user of the failure prediction although the two HDDs are unlikely to fail at the same time. The user, however, does not know what to do even when he/she is notified of the failure prediction.\nIn relation to the above described technique, there has been disclosed a technique that S.M.A.R.T. information is stored in a server in case that the S.M.A.R.T. information is corrupted, and in the server, the S.M.A.R.T. information is compared with a previous log to thus reduce the load for analyzing the cause of a failure (see, for example, Japanese Laid Open Patent Publication (Kokai) No. 2001-312375).\nAlso, there has been disclosed a technique that a failure is predicted based on, for example, the number of accesses to an HDD or the like (see, for example, Japanese Laid Open Patent Publication (Kokai) No. 2006-256251).\nFor example, according to the technique described in Japanese Laid Open Patent Publication (Kokai) No. 2001-312375, a server can be notified of a status of an image forming apparatus, and also, values obtained using the failure predicting function can be backed up to the server.\nAccording to the technique described in Japanese Laid Open Patent Publication (Kokai) No. 2006-256251, a failure can be predicted based on the number of accesses. These conventional arts can improve the accuracy of the failure predicting function.\nHowever, when both the mirroring function and the failure predicting function are enabled, problems described hereafter arise. First, unnecessary notification of a failure prediction is provided to a user. This means that, for example, an image forming apparatus provides notification of a failure prediction to a user interface or the like.\nNext, unnecessary notification of a failure prediction is provided to a service engineer. This means that, for example, unnecessary notification of an error is provided to a server that manages an image forming apparatus, and a service engineer is called.\nMoreover, notification of a failure prediction is provided which causes a service engineer to falsely recognize that a trouble is serious and urgent. This means that, for example, notification of an error is provided to a server that manages an image forming apparatus, and the service engineer is called although a visit of the service engineer is unnecessary.\nFurther, even if a display indicative of a failure prediction notification is produced on a user interface of an image forming apparatus or the like, a user does not know what to do.\nAs described above, the conventional techniques have the problem of requiring unnecessary operational costs for calling a service engineer, replacing an HDD, and so on."}
-{"text": "Lithium ion secondary batteries are widely used for portable electronic instruments such as mobile phones or notebook-size personal computers. As a cathode active material for a lithium ion secondary battery, a composite oxide of lithium with a transition metal, etc. (hereinafter sometimes referred to as lithium-containing composite oxide) such as LiCoO2, LiNiO2, LiNi0.8Co0.2 O2 or LiMn2O4, is employed. In recent years, it is desired to further reduce the size and weight as a lithium ion secondary battery for portable electronic instruments or vehicles, and a further improvement in the discharge capacity per unit mass and in the characteristics such that the discharge capacity will not be decreased after charge and discharge cycles are repeatedly carried out (hereinafter sometimes referred to as cycle characteristic), is desired.\nTo improve the cycle characteristics, it has been known to be effective to form a covering layer on a lithium-containing composite oxide.\nPatent Document 1 discloses a process of adding an ammonium fluoride aqueous solution to a dispersion having a lithium-containing composite oxide dispersed in an aluminum nitrate aqueous solution, followed by filtration, washing and heating to form a covering layer of aluminum fluoride on the surface of the lithium-containing composite oxide. However, by this process, a step of carrying out filtration and washing is essential, whereby the process is complicated and in addition, a waste disposal treatment is necessary, and thus the process is inferior in the productivity. Further, when the wet cake after filtration is dried, the cathode active material tends to be agglomerated to form coarse particles.\nPatent Document 2 discloses a method of fluorinating a cathode material using a fluorinating agent such as fluorine (F2), nitrogen trifluoride or chlorine trifluoride. However, this method requires special equipment since a highly toxic gas is used.\nPatent Document 3 discloses a process of stirring and mixing an aqueous solution containing zirconium and a lithium-containing composite oxide represented by the formula LipNxMyOzFa (0.9\u2266p\u22661.1, 0.965\u2266x<1.00, 0<y\u22660.035, 1.9\u2266z\u22662.1, x+y=1 and 0\u2266a\u22660.02) wherein the molar amount of Li element is from 0.9 to 1.1 molar times the total molar amount of the transition metal element, and firing the mixture at high temperature of at least 450\u00b0 C. in an oxygen atmosphere to obtain a cathode active material having the surface layer of the lithium-containing composite oxide covered with zirconium oxide. By this method, it is difficult to cover the lithium-containing composite oxide with a compound other than an oxide."}
-{"text": "The invention relates to novel phenyl-substituted cyclic ketoenols, to a plurality of processes and intermediates for their preparation and to their use as pesticides.\nIt has akready been discosed that certain phenyl-substituted cyc ketoenols are active as insectidces, acaricides and/or herbicides.\nAlso known are 1H-arylpyrrolidine-dione derivatives (BP-A-456 063, EP-A-521 334, EP-A-596 298, EP-A-613 884, EP-A-613 885, DE 44 40 594, WO 94/01 997, WO 95/01 358, WO 95126 954, WO 95/20 572, EP-A-0 668 267, WO 96/15 395, WO 96/35 664, WO 97/01 535 and WO 97/02 243).\nIt is known that certain substituted xcex943-dihydrofuran-2-one derivatives have herbicidal properties (cf. DE-A-4 014 420). The synthesis of the tetronic acid derivatives used as staring materials (such as, for example, 3-(2-methyl-phenyl)-4-hydroxy-5-(4-fluorophenyl)-xcex943-dihydrofuran-2-one) is likewise described in DE-A-4 014 420. Compounds of similar struftre without details of an insecticidal and/or acaricidal activity are known from the publication Campbell et al., J. Chem. Soc., Perkin Trans. 1, 1985, (S) 1567-76. Furthermore, 3-aryl-xcex943-dihydrofuranone derivatives having herbicidal, acaricidal and insecticidal properties are disclosed in EP-A-528 156, EP-A 0 647 637, WO 95/26345, WO 96/20 196, WO 96/25 395, WO 96/35 664, WO 97/01 535 and WO 97/02 243.\nCertain phenytl-pyrone derivatives unsubstituted in the phenyl ring have already been dislosed (cf. AM. Chirazi, T. Kappe and E. Ziegler, Arch. Phatin. 309, 558 (1976) and K. -H. Boltze and K. Heidenbluth, Chem. Ber. 91, 2849), a possible utility for these compounds as pesticides not being indicated. Phenyl-pyrone dertves substituted in the phenyl ring and having herbicidal, acaricidal and insecticidal propeties are descnrbed in EP-A-588 137, WO 96/25 395, WO 96/35 664 WO 97/01 535 and WO 97/02 243.\nHowever, the acaricidal and insecticidal activity and/or spectrum of action, and/or the toleration of the known compounds by plants, in parficular by crops, is not always satisfactory.\nThe invention, accordingly, provides compounds of the formula (I) \nin which\nX represents halogen, alkyl, alkenyl, alkynyl, alkoxy, benzyloxy, halogenoalkyl, halogenoalkoxy, cyano or nitro,\nZ represents hydrogen, amino, halogen, alkyl, alkoxy, halogenoalkyl, halogenoalkoxy, hydroxyl, cyano, nitro or respectively optionally substituted phenoxy, phenylthio, 5- or 6-membered hetaryloxy, 5- or 6-membered hetarylthio, phenylalkyloxy or phenylalkylthio and\nHet represents one of the groups \nxe2x80x83in which\nA represents a respectively optionally substituted radical from the group consisting of alkyl, alkenyl, alkoxyalkyl, polyalkoxyalkyl and alkylthioalkyl, represents respectively saturated or unsaturated and optionally substituted cycloalkyl or heterocyclyl or represents respectively optionally halogen-, alkyl-, halogenoalkyl-, alkoxy-, hatogenoalkoxy-, cyano- or nitro-substituted aryl, arylalkyl or hetaryl,\nB represents alkyl or alkoxyalkyl or\nA and B together with the carbon atom that they are attached to represent a saturated or unsaturated, optionally substituted carbocycle or heterocycle,\nD represents hydrogen or represents an optionally substituted radical from the group consisting of alkyl, alkenyl, alkynyl, alkoxyalkyl, polyalkoxyalkyl, alkylthioalkyl, saturated or unsaturated cycloalkyl, saturated or unsaturated heterocyclyl, arylalkyl, aryl, hetarylalkyl or hetaryl or\nA and D together with the atoms that they are attached to represent a respectively optionally substituted carbocycle or heterocycle,\nG represents hydrogen (a) or represents one of the groups \nxe2x80x83in which\nE represents a metal ion equivalent or an ammonium ion,\nL represents oxygen or sulphur,\nM represents oxygen or sulphur,\nR1 represents respectively optionally halogen-substituted alkyl, alkenyl, alkoxyalkyl, alkylthioalkyl or polyalkoxyalkyl or represents respectively optionally halogen-, alkyl- or alkoxy-substituted cycloalkyl or heterocyclyl or represents respectively optionally substituted phenyl, phenylalkyl, hetaryl, phenoxyalkyl or hetarytoxyalkyl,\nR2 represents respectively optionally halogen-substituted alkyl, alkenyl, alkoxyalkyl or polyalkoxyalkyl or represents respectively optionally substituted cycloalkyl, phenyl or benzyl,\nR3, R4 and R5 independently of one another each represent respectively optionally halogen-substituted alkyl, alkoxy, alkylamino, dialkylamino, alkylthio, alkenylthio or cycloalkylthio or represent respectively optionally substituted phenyl, benzyl, phenoxy or phenylthio,\nR6 and R7 independently of one another each represent hydrogen, represent respectively optionally halogen-substituted alkyl, cycloalkyl, alkenyl, alkoxy, alkoxyalkyl, represent respectively optionally substituted phenyl or benzyl, or form together with the nitrogen atom that they are attached to an optionally oxygen- or sulphur-containing, optionally substituted cycle.\nThe compounds of the formula (I) can also be present, depending on the nature of the substituents, as geometric and/or optical isomers and isomer mixtures of differing composition which, if appropriate, can be separated in a customary manner. Both the pure isomers and the isomer mixtures, their preparation and use, and compositions comprising them are part of the subject matter of the present invention. In the following, for simplicity, however, compounds of the formula (I) are always referred to, although both pure compounds and, if appropriate, mixtures having different proportions of isomer compounds are intended.\nIncluding the meanings (1) to (4) of the group Het, the following principal structures (I-1) to (I-4) result: \nin which\nA, B, D, G, X and Z are each as defined above.\nIncluding the various meanings (a), (b), (c), (d), (e), (f) and (g) of the group G, the following principal structures (I-1-a) to (I-1-g) result if Het represents the group (1) \nin which\nA, B, E, L, M, X, Z, R1, R2, R3, R4, R5, R6 and R7 are each as defined above.\nIncluding the various meanings (a), (b), (c), (d), (e), (f) and (g) of the group G, the following principal structures (I-2-a) to (I-2-g) result if Het represents the group (2) \nin which\nA, B, E, L, M, X, Z, R1, R2, R3, R4, R5, R6 and R7 are each as defined above.\nIncluding the various meanings (a), (b), (c), (d), (e), (f) and (g) of the group G, the following principal structures (I-3-a) to (I-3-g) result if Het represents the group (3) \nin which\nA, B, E, L, M, X, Z, R1, R2 R3, R4, R5, R6 and R7 are each as defined above.\nDepending on the position of the substituent G, the compounds of the formula (I-4) can be present in the two isomeric forms of formulae (I-4)a and (I-4)b\nwhich is intended to be expressed by the dashed line in the formula (I-4).\nThe compounds of the formulae (I-4)a and (I-4)b can be present both as mixtures and in the form of their pure isomers. Mixtures of the compounds of the formulae (I4)a and (I-4)b can, if desired, be separated by physical methods in a manner known per se, for example by chromatographic methods.\nFor better clarity, in the following in each case only one of the possible isomers is shown. This does not exclude the possibility that the compounds can optionally be present in the form of the isomer mixtures or in the other respective isomer form.\nIncluding the various meanings (a), (b), (c), (d), (e), (f) and (g) of the group G, the following principal structures (I-4-a) to (I-4-g) result if Het represents the group (4) \nin which\nA, D, E, L, M, X, Z, R1, R2, R3, R4, R5, R6 and R7 are each as defined above.\nFurthermore, it has been found that the novel compounds of the formula (I) are 10 obtained by one of the processes described below:\n(A) Compounds of the formula (I-1-a) \nin which\nA, B, X and Z are each as defined above, are obtained by the intramolecular condensation of compounds of the formula (II) \nxe2x80x83in which\nA, B, X and Z are each as defined above, and\nR8 represents alkyl (preferably C1-C6-alkyl), in the presence of a diluent and in the presence of a base.\n(B) Furthermore, it was found that compounds of the formula (I-2-a) \nin which\nA, B, X and Z are each as defined above,\nare obtained by the intramolecular condensation of compounds of the formula (III) \nxe2x80x83in which\nA, B, X, Z and R8 are each as defined above,\nin the presence of a diluent and in the presence of a base.\n(C) Furthermore, it was found that compounds of the formula (I-3-a) \nin which\nA, B, X and Z are each as defined above,\nare obtained by the intramolecular cyclization of compounds of the formula (IV) \nxe2x80x83in which\nA, B, X, Z and R8 are each as defined above and\nW represents hydrogen, halogen, alkyl (preferably C1-C6-alkyl) or alkoxy (preferably C1-C8-alkoxy), if appropriate in the presence of a diluent and in the presence of an acid.\n(D) Furthermore, it was found that compounds of the formula (I-4-a) \nin which\nA, D, X and Z are each as defined above,\nare obtained by reacting compounds of the formula (V) \nxe2x80x83in which\nA and D are each as defined above,\nor their silyl enol ethers of the formula (Va) \nxe2x80x83in which\nA and D are each as defined above and\nR8xe2x80x2 represents alkyl (preferably methyl),\nwith compounds of the formula (VI) \nxe2x80x83in which\nX and Z are each as defined above and\nHal represents halogen (preferably chlorine or bromine),\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid acceptor.\nFurthermore, it was found\n(E) that the compounds of the formulae (I-1-b) to (I-4-b) shown above in which A, B, D, R1, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above\nxcex1) with acid halides of the formula (VII) \nxe2x80x83in which\nR1 is as defined above and\nHal represents halogen (in particular chlorine or bromine), or\nxcex2) with carboxylic anhydrides of the formula (VIII)\nR1xe2x80x94COxe2x80x94Oxe2x80x94COxe2x80x94R1xe2x80x83xe2x80x83(VIII)\nxe2x80x83in which\nR1 is as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent;\n(F) that the compounds of the formulae (I-1-c) to (I-4-c) shown above in which A, B, D, R2, M, X and Z are each as defined above and L represents oxygen are obtained by reacting compounds of the formulae (I-1-a) to\n(I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nwith chloroformic esters or chloroformic thioesters of the formula (IX)\nR2xe2x80x94Mxe2x80x94COxe2x80x94Clxe2x80x83xe2x80x83(IX)\nxe2x80x83in which\nR2 and M are each as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent;\n(G) that compounds of the formulae (I-1-c) to (I-4-c) shown above in which A, B, D, R2, M, X and Z are each as defined above and L represents sulphur are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nwith chloromonothioformic esters or chlorodithioformic esters of the formula (X) \nxe2x80x83in which\nM and R2 are each as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-bindina agent,\n(H) that compounds of the formulae (I-1-d) to (I-4-d) shown above in which A, B, D, R3, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nwith sulphonyl chlorides of the formula (XII)\nR3xe2x80x94SO2xe2x80x94Clxe2x80x83xe2x80x83(XII)\nxe2x80x83in which\nR3 is as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent,\n(I) that compounds of the formulae (I-1-e) to (I-4-e) shown above in which A, B, D, L, R4, R5, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nwith phosphorus compounds of the formula (XIII) \nxe2x80x83in which\nL, R4 and R5 are each as defined above and\nHal represents halogen (in particular chlorine or bromine),\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent,\n(J) that compounds of the formulae (I-1-f) to (I-4-f) shown above in which A, B, D, E, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) in which A, B, D, X and Z are each as defined above, in each case\nwith metal compounds or amines of the formulae (XIV) or (XV) \nxe2x80x83in which\nMe represents a mono- or divalent metal (preferably an alkali metal or alkaline earth metal, such as lithium, sodium, potassium, magnesium or calcium),\nt represents the number 1 or 2 and\nR10, R11, R12 independently of one another each represent hydrogen or alkyl (preferably C1-C8-alkyl),\nif appropriate in the presence of a diluent;\n(K) that compounds of the formulae (I-1-g) to (I-4-g) shown above in which A, B, D, L, R6, R7, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nxcex1) with isocyanates or isothiocyanates of the formula (XVI)\nR6xe2x80x94Nxe2x95x90Cxe2x95x90Lxe2x80x83xe2x80x83(XVI)\nxe2x80x83in which\nR6 and L are each as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of a catalyst, or\nxcex2) with carbamoyl chlorides or thiocarbamoyl chlorides of the formula (XVII) \nxe2x80x83in which\nL, R6 and R7 are each as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nFurthermore, it has been found that the novel compounds of the formula (I) have a very good activity as pesticides, preferably as insecticides and acaricides, and that they additionally are very well tolerated by plants, in particular by crops.\nFormula (I) provides a general definition of the compounds according to the invention. Preferred substituents or ranges of the radicals shown in the formulae mentioned hereinabove and hereinbelow are illustrated below:\nX preferably represents halogen, C1-C6-alkyl, C2-C6-alkenyl, C2-C6-alkynyl, C1-C6-alkoxy, benzyloxy, C1-C4-halogenoalkyl, C1-C4-halogenoalkoxy, cyano or nitro.\nZ preferably represents hydrogen, amino, halogen, C1-C6-alkyl, C1-C6-alkoxy, C1-C4-halogenoalkyl, C1-C4-halogenoalkoxy, hydroxyl, cyano, nitro or respectively optionally halogen-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C4-halogenoalkyl-, C1-C4-halogenoalkoxy-, nitro- or cyano-substituted phenoxy, phenylthio, thiazolyloxy, pyridinyloxy, pyrimidinyloxy, pyrazolyloxy, phenyl-C1-C4-alkyloxy or phenyl-C1-C7-alkylthio.\nHet preferably represents one of the groups \nA preferably represents respectively optionally halogen-substituted C1-C12-alkyl, C2-C8-alkenyl, C1-C10-alkoxy-C1-C8-alkyl, poly-C1-C8-alkoxy-C1-C8-alkyl or C1-C10-alkylthio-C1-C6-alkyl, preferably represents optionally halogen-, C1-C6-alkyl- or C1-C6-alkoxy-substituted C3-C8-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or preferably represents respectively optionally halogen-, C1-C6-alkyl-, C1-C6-halogenoalkyl-, C1-C6-alkoxy-, C1-C6-halogenoalkoxy-, cyano- or nitro-substituted phenyl, naphthyl, phenyl-C1-C6-alkyl, naphthyl-C1-C6-alkyl or hetaryl having 5 or 6 ring atoms and one to three hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of furanyl, pyridyl, imidazolyl, triazolyl, pyrazolyl, indolyl, thiazolyl and thienyl).\nB preferably represents C1-C12-alkyl or C1-C8-alkoxy-C1-C6-alkyl or\nA, B and the carbon atom that they are attached to preferably represent C3-C10-cycloalkyl or C1-C10-cycloalkenyl where in each case one methylene group is optionally replaced by oxygen or sulphur and which are optionally substituted by C1-C8-alkyl, C3-C10-cycloalkyl, C1-C8-halogenoalkyl, C1-C8-alkoxy, C1-C8-alkylthio, halogen or phenyl or\nA, B and the carbon atom that they are attached to preferably represent C5-C6-cycloalkyl which is substituted by an alkylenediyl group optionally containing one or two not directly adjacent oxygen and/or sulphur atoms, or by an alkylenedioxy group or an alkylenedithioyl group forming a further five- to eight-membered ring with the carbon atom that it is attached to, or\nA, B and the carbon atom that they are attached to preferably represent C3-C8-cycloalkyl or C5-C8-cycloalkenyl in which two carbon atoms are linked to each other by respectively optionally C1-C6-alkyl-, C1-C6-alkoxy- or halogen-substituted C3-C6-alkanediyl, C1-C6-alkenediyl or C4-C6-alkanedienediyl, one methylene group in each case being optionally replaced by oxygen or sulphur.\nD preferably represents hydrogen, preferably represents respectively optionally halogen-substituted C1-C12-alkyl, C3-C8-alkenyl, C3-C8-alkynyl, C1-C10-alkoxy-C2-C8-alkyl, poly-C1-C8-alkoxy-C2-C8-alkyl or C1-C10-alkylthio-C2-C8-alkyl, preferably represents optionally halogen-, C1-C4-alkyl-, C1-C4-alkoxy- or C1-C4-halogenoalkyl-substituted C3-C8-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or preferably represents respectively optionally halogen-, C1-C6 alkyl-, C1-C6-halogenoalkyl-, C1-C6-alkoxy-, C1-C6-halogenoalkoxy-, cyano- or nitro-substituted phenyl, hetaryl having 5 or 6 ring atoms and one or two hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of furanyl, imidazolyl, pyridyl, thiazolyl, pyrazolyl, pyrimidyl, pyrrolyl, thienyl and triazolyl), phenyl-C1-C6-alkyl or hetaryl-C1-C6-alkyl having 5 or 6 ring atoms and one or two hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of furanyl-, imidazolyl-, pyridyl-, thiazolyl-, pyrazolyl-, pyrimidyl-, pyrrolyl-, thienyl- and triazolyl-C1-C6-alkyl) or\nA and D together preferably represent a C3-C6-alkanediyl, C3-C6-alkenediyl or C4-C6-alkadienediyl group in which respectively optionally one methylene group is replaced by oxygen or sulphur and which are respectively optionally substituted by halogen or respectively optionally halogen-substituted C1-C10-alkyl, C1-C6-alkoxy, C1-C6-alkylthio, C3-C7-cycloalkyl, phenyl or benzyloxy or by a further C3-C6-alkanediyl, C3-C6-alkenediyl or C4-C6-alkadienediyl group forming a fused ring, in which optionally respectively one methylene group is replaced by oxygen or sulphur and which are optionally substituted by C1-C6-alkyl, or\nA and D together represent a C3-C6-alkanediyl or C3-C6-alkenediyl group containing in each case optionally one of the following groups \nG preferably represents hydrogen (a) or represents one of the groups \nxe2x80x83in which\nE represents a metal ion equivalent or an ammonium ion,\nL represents oxygen or sulphur and\nM represents oxygen or sulphur.\nR1 preferably represents respectively optionally halogen-substituted C1-C20-alkyl, C2-C20-alkenyl, C1-C8-alkoxy-C1-C8-alkyl, C1-C8-alkylthio-C1-C8-alkyl or poly-C1-C8-alkoxy-C1-C8-alkyl or preferably represents optionally halogen-, C1-C6-alkyl- or C1-C6-alkoxy-substituted C3-C8-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur,\npreferably represents optionally halogen-, cyano-, nitro-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C6-halogenoalkyl-, C1-C6-halogenoalkoxy-, C1-C6-alkylthio- or C1-C6-alkylsulphonyl-substituted phenyl,\npreferably represents optionally halogen-, nitro-, cyano-, C1-C6-alkyl, C1-C6-alkoxy-, C1-C6-halogenoalkyl- or C1-C6-halogenoalkoxy-substituted phenyl-C1-C6-alkyl,\npreferably represents optionally halogen- or C1-C6-alkyl-substituted 5- or 6-membered hetaryl having one or two hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of pyrazolyl, thiazolyl, pyridyl, pyrimidyl, furanyl and thienyl), preferably represents optionally halogen- or C1-C6-alkyl-substituted phenoxy-C1-C6-alkyl or\npreferably represents optionally halogen-, amino- or C1-C6-alkyl-substituted 5- or 6-membered hetaryloxy-C1-C6-alkyl having one or two hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of pyridyloxy-C1-C6-alkyl, pyrimidyloxy-C1-C6-alkyl and thiazolyloxy-Ci-C6-alkyl).\nR2 preferably represents respectively optionally halogen-substituted C1-C20-alkyl, C2-C20-alkenyl, C1-C8-alkoxy-C2-C8-alkyl or poly-C1-C8-alkoxy-C2-C8-alkyl,\npreferably represents optionally halogen-, C1-C6-alkyl- or C1-C6-alkoxy-substituted C3-C8-cycloalkyl or\npreferably represents respectively optionally halogen-, cyano-, nitro-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C6-halogenoalkyl- or C1-C6-halogenoalkoxy-substituted phenyl or benzyl.\nR3 preferably represents optionally halogen-substituted C1-C8-alkyl or respectively optionally halogen-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C4-halogenoalkyl-, C1-C4-halogenoalkoxy-, cyano- or nitro-substituted phenyl or benzyl.\nR4 and R5 independently of one another each preferably represent respectively optionally halogen-substituted C1-C8-alkyl, C1-C8-alkoxy, C1-C8-alkylamino, di-(C1-C8-alkyl)-amino, C1-C8-alkylthio or C3-C8-alkenylthio or preferably represent respectively optionally halogen-, nitro-, cyano-, C1-C4-alkoxy-, C1-C4-halogenoalkoxy-, C1-C4-alkylthio-, C1-C4-halogenoalkylthio-, C1-C4-alkyl- or C1-C4-halogenoalkyl-substituted phenyl, phenoxy or phenylthio.\nR6 and R7 independently of one another each preferably represent hydrogen, preferably represent respectively optionally halogen-substituted C1-C8-alkyl, C3-C8-cycloalkyl, C1-C8-alkoxy, C3-C8-alkenyl or C1-C8-alkoxy-C2-C8-alkyl, preferably represent respectively optionally halogen-, C1-C8-alkyl-, C1-C8-halogenoalkyl- or C1-C8-alkoxy-substituted phenyl or benzyl or together preferably represent an optionally C1-C6-alkyl-substituted C3-C6-alkylene radical in which optionally one methylene group is replaced by oxygen or sulphur.\nR13 preferably represents hydrogen or respectively optionally halogen-substituted C1-C8-alkyl or C1-C8-alkoxy, preferably represents optionally halogen-, C1-C4-alkyl- or C1-C4-alkoxy-substituted C3-C8-cycloalkyl in which optionally one methylene group is replaced by oxygen or sulphur, or preferably represents respectively optionally halogen-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C4-halogenoalkyl-, C1-C4-halogenoalkoxy-, nitro- or cyano-substituted phenyl, phenyl-C1-C4-alkyl or phenyl-C1-C4-alkoxy.\nR14 preferably represents hydrogen or C1-C8-alkyl or\nR13 and R14 together preferably represent C4-C6-alkanediyl.\nR15 and R16 are identical or different and each preferably represent C1-C6-alkyl or\nR15 and R16 together preferably represent a C2-C4-alkanediyl radical which is optionally substituted by C1-C6-alkyl or by optionally halogen-, C1-C4-alkyl-, C1-C4-halogenoalkyl-, C1-C4-alkoxy-, C1-C4-halogenoalkoxy-, nitro- or cyano-substituted phenyl.\nR17 and R18 independently of one another each preferably represent hydrogen, preferably represent optionally halogen-substituted C1-C8-alkyl or preferably represent optionally halogen-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C4-halogenoalkyl-, C1-C4-halogenoalkoxy-, nitro- or cyano-substituted phenyl or\nR17 and R18 together with the carbon atom that they are attached to preferably represent optionally C1-C4-alkyl-substituted C5-C7-cycloalkyl in which optionally one methylene group is replaced by oxygen or sulphur.\nR19 and R20 independently of one another each preferably represent C1-C10-alkyl, C1-C10-alkenyl, C1-C10-alkoxy, C1-C10-alkylamino, C3-C10-alkenylamino, di-(C1-C10-alkyl)-amino or di-(C3-C10-alkenyl)-amino.\nX particularly preferably represents fluorine, chlorine, bromine, C1-C4-alkyl, C1-C4-alkoxy, benzyloxy, C1-C2-halogenoalkyl, C1-C2-halogenoalkoxy, cyano or nitro.\nZ particularly preferably represents hydrogen, amino, fluorine, chlorine, bromine, C1-C4-alkyl, C1-C4-alkoxy, C1-C2-halogenoalkyl, C1-C2-halogenoalkoxy, hydroxyl, cyano, nitro or respectively optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C2-halogenoalkyl-, C1-C2-halogenoalkoxy-, nitro- or cyano-substituted phenoxy or benzyloxy.\nHet particularly preferably represents one of the groups \nA particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C10-alkyl, C2-C6-alkenyl, C1-C8-alkoxy-C1-C6-alkyl, poly-C1-C6-alkoxy-C1-C6-alkyl or C1-C8-alkylthio-C1-C6-alkyl or particularly preferably represents optionally fluorine-, chlorine-, C1-C4-alkyl- or C1-C4-alkoxy-substituted C3-C7-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-halogenoalkyl-, C1-C4-alkoxy-, C1-C4-halogenoalkoxy-, cyano, or nitro-substituted phenyl, furanyl, pyridyl, imidazolyl, triazolyl, pyrazolyl, indolyl, thiazolyl, thienyl or phenyl-C1-C4-alkyl.\nB particularly preferably represents C1-C10-alkyl or C1-C6-alkoxy-C1-C4-alkyl or\nA, B and the carbon atom that they are attached to particularly preferably represent C3-C8-cycloalkyl or C5-C8-cycloalkenyl where one methylene group in each case is optionally replaced by oxygen or sulphur and which are optionally substituted by C1-C6-alkyl, C3-C8-cycloalkyl, C1-C3-haloogenoalkyl, C1-C6-alkoxy, C1-C6-alkylthio, fluorine, chlorine or phenyl or\nA, B and the carbon atom that they are attached to particularly preferably represent C1-C6-cycloalkyl which is substituted by an alkylenediyl group optionally containing one or two not directly adjacent oxygen or sulphur atoms or by an alkylenedioxy group or by an alkylenedithiol group which forms together with the carbon atom that it is attached to a further five- to seven-membered ring or\nA, B and the carbon atom that they are attached to particularly preferably represent C3-C6-cycloalkyl or C5-C6-cycloalkenyl in which two carbon atoms are linked to each other by respectively optionally C1-C4-alkyl-, C1-C4-alkoxy-, fluorine-, chlorine- or bromine-substituted C3-C5-alkanediyl or C3-C5-alkenediyl, one methylene group in each case being optionally replaced by oxygen or sulphur, or are linked to each other by butadienediyl.\nD particularly preferably represents hydrogen, particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl or C1-C6-alkoxy, particularly preferably represents optionally fluorine-, chlorine-, C1-C4-alkyl-, C1-C4-alkoxy- or C1-C2-halogenoalkyl-substituted C3-C7-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-halogenoalkyi, C1-C4-alkoxy-, C1-C4-halogenoalkoxy-, cyano- or nitro-substituted phenyl, furanyl, imidazolyl, pyridyl, thiazolyl, pyrazolyl, pyrimidyl, pyrrolyl, thienyl, triazolyl or phenyl-C1-C4-alkyl or\nA and D together particularly preferably represent a C3-C5-alkanediyl or C1-C5-alkenediyl group in which in each case one methylene group is optionally replaced by oxygen or sulphur and which are optionally substituted by fluorine, chlorine or by respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl, C1-C4-alkoxy, C1-C4-alkylthio, C3-C6-cycloalkyl, phenyl or benzyloxy or\nin which in each case optionally one of the following groups is contained: \nG particularly preferably represents hydrogen (a) or represents one of the groups \nxe2x80x83in which\nE represents a metal ion equivalent or an ammonium ion,\nL represents oxygen or sulphur and\nM represents oxygen or sulphur.\nR1 particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C 16-alkyl, C2-C16-alkenyl, C1-C6-alkoxy-C1-C6-alkyl, C1-C6-alkylthio-C1-C6-alkyl or poly-C1-C6-alkoxy-C1-C6-alkyl or particularly preferably represents optionally fluorine-, chlorine-, C1-C5-alkyl- or C1-C5-alkoxy-substituted C3-C7-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur,\nparticularly preferably represents optionally fluorine-, chlorine-, bromine-, cyano-, nitro-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C3-halogenoalkyl-, C1-C3-halogenoalkoxy-, C1-C4-alkylthio- or C1-C4-alkylsulphonyl-substituted phenyl,\nparticularly preferably represents optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C3-halogenoalkyl- or C1-C3-halogeno-alkoxy-substituted phenyl-C1-C4-alkyl,\nparticularly preferably represents respectively optionally fluorine-, chlorine-, bromine- or C1-C4-alkyl-substituted pyrazolyl, thiazolyl, pyridyl, pyrimidyl, furanyl or thienyl,\nparticularly preferably represents optionally fluorine-, chlorine-, bromine- or C1-C4-alkyl-substituted phenoxy-C1-C5-alkyl or\nparticularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, amino- or C1-C4-alkyl-substituted pyridyloxy-C1-C5-alkyl, pyrimidyloxy-C1-C5-alkyl or thiazolyloxy-C1-C5-alkyl.\nR2 particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C16-alkyl, C2-C16-alkenyl, C1-C6-alkoxy-C2-C6-alkyl or poly-C1-C6-alkoxy-C2-C6-alkyl,\nparticularly preferably represents optionally fluorine-, chlorine-, C1-C4-alkyl- or C1-C4-alkoxy-substituted C3-C7-cycloalkyl or\nparticularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, cyano-, nitro-, C1-C4-alkyl-, C1-C3-alkoxy-, C1-C3-halogenoalkyl- or C1-C3-halogenoalkoxy-substituted phenyl or benzyl.\nR3 particularly preferably represents optionally fluorine- or chlorine-substituted C1-C6-alkyl or respectively optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C2-halogenoalkoxy-, C1-C2-halogenoalkyl-, cyano-, or nitro-substituted phenyl or benzyl.\nR4 and R5 independently of one another each particularly preferably represent respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl, C1-C6-alkoxy, C1-C6-alkylamino, di-(C1-C6-alkyl)-amino, C1-C6-alkylthio or C3-C4-alkenylthio or particularly preferably represent respectively optionally fluorine-, chlorine-, bromine-, nitro-, cyano-, C1-C3-alkoxy-, C1-C3-haloge-noalkoxy-, C1-C3-alkylthio-, C1-C3-halogenoalkylthio-, C1-C3-alkyl- or C1-C3-halogenoalkyl-substituted phenyl, phenoxy or phenylthio.\nR6 and R7 independently of one another each particularly preferably represent hydrogen, particularly preferably represent respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl, C3-C6-cycloalkyl, C1-C6-alkoxy, C3-C6-alkenyl or C1-C6-alkoxy-C2-C6-alkyl, particularly preferably represent respectively optionally fluorine-, chlorine-, bromine-, C1-C5-halogenoalkyl-, C1-C5-alkyl- or C1-C5-alkoxy-substituted phenyl or benzyl, or together particularly preferably represent an optionally C1-C4-alkyl-substituted C3-C6-alkylene radical in which optionally one methylene group is replaced by oxygen or sulphur.\nR13 particularly preferably represents hydrogen or respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl or C1-C6-alkoxy, particularly preferably represents optionally fluorine-, C1-C2-alkyl- or C1-C2-alkoxy-substituted C3-C7-cycloalkyl in which optionally one methylene group is replaced by oxygen or sulphur, or particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, C1-C5-alkyl-, C1-C5-alkoxy-, C1-C2-halogenoalkyl-, C1-C2-halogenoalkoxy-, nitro- or cyano-substituted phenyl, phenyl-C1-C3-alkyl or phenyl-C1-C2-alkyloxy.\nR14 particularly preferably represents hydrogen or C1-C6-alkyl or\nR13 and R14 together particularly preferably represent C4-C6-alkanediyl.\nR15 and R16 are identical or different and each particularly preferably represent C1-C4-alkyl or\nR15 and R16 together particularly preferably represent a C2-C3-alkanediyl radical which is optionally substituted by C1-C4-alkyl or by optionally fluorine-, chlorine-, bromine-, C1-C2-alkyl-, C1-C2-halogenoalkyl-, C1-C2-alkoxy-, C1-C2-halogenoalkoxy-, nitro- or cyano-substituted phenyl.\nX very particularly preferably represents fluorine, chlorine, bromine, methyl, ethyl, n-propyl, iso-propyl, n-butyl, iso-butyl, tert-butyl, methoxy, ethoxy, n-propoxy, iso-propoxy, trifluoromethyl, trifluoromethoxy, difluoromethoxy, cyano or nitro,\nZ very particularly preferably represents hydrogen, amino, fluorine, chlorine, bromine, methyl, ethyl, n-propyl, iso-propyl, n-butyl, iso-butyl, tert-butyl, methoxy, ethoxy, n-propoxy, iso-propoxy, trifluoromethyl, trifluoromethoxy, difluoromethoxy, cyano or nitro.\nHet very particularly preferably represents one of the groups \nA very particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C8-alkyl, C2-C4-alkenyl, C1-C6-alkoxy-C1-C4-alkyl, poly-C1-C4-alkoxy-C1-C4-alkyl or C1-C6-alkylthio-C1-C4-alkyl, or very particularly preferably represents optionally fluorine-, chlorine-, methyl- or methoxy-substituted C3-C6-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or very particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, methyl-, ethyl-, n-propyl-, iso-propyl-, methoxy-, ethoxy-, trifluoromethyl-, trifluoromethoxy-, cyano- or nitro-substituted phenyl, pyridyl or benzyl.\nB very particularly preferably represents C1-C8-alkyl or C1-C4-alkoxy-C1-C2-alkyl or\nA, B and the carbon atom that they are attached to very particularly preferably represent C3-C8-cycloalkyl or C5-C8-cycloalkenyl in which in each case optionally one methylene group is replaced by oxygen or sulphur and which are optionally substituted by methyl, ethyl, n-propyl, iso-propyl, butyl, iso-butyl, sec-butyl, tert-butyl, cyclohexyl, trifluoromethyl, methoxy, ethoxy, n-propoxy, iso-propoxy, butoxy, iso-butoxy, sec-butoxy, tert-butoxy, methylthio, ethylthio, fluorine, chlorine or phenyl or\nA, B and the carbon atom that they are attached to very particularly preferably represent C3-C6-cycloalkyl or C5-C6-cycloalkenyl in which two carbon atoms are linked together by C3-C4-alkanediyl or C3-C4-alkenediyl in which in each case optionally one methylene group is replaced by oxygen or sulphur, or are linked together by butadienediyl.\nD very particularly preferably represents hydrogen, very particularly preferably represents respectively optionally fluorine-substituted C1-C4-alkyl or C1-C4-alkoxy or C3-C6-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or very particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, methyl-, ethyl-, n-propyl-, iso-propyl-, methoxy-, ethoxy-, trifluoromethyl-, trifluoromethoxy-, cyano- or nitro- substituted phenyl, furanyl, pyridyl, thienyl or benzyl, or\nA and D together very particularly preferably represent a C3-C5-alkanediyl or C3-C5-alkenediyl group in which in each case optionally one methylene group is replaced by oxygen or sulphur and which are optionally substituted by fluorine, chlorine or by respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl, C1-C4-alkoxy, C1-C4-alkylthio, C3-C6-cycloalkyl, phenyl or benzyloxy.\nG very particularly preferably represents hydrogen (a) or represents one of the groups \nxe2x80x83in which\nE represents a metal ion equivalent or an ammonium ion,\nL represents oxygen or sulphur and\nM represents oxygen or sulphur.\nR1 very particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C14-alkyl, C2-C14-alkenyl, C1-C4-alkoxy-C1-C6-alkyl, C1-C4-alkylthio-C1-C6-alkyl, poly-C1-C4-alkoxy-C1-C4-alkyl or very particularly preferably represents optionally fluorine-, chlorine-, methyl-, ethyl-, n-propyl-, i-propyl-, n-butyl-, i-butyl-, tert-butyl-, methoxy-, ethoxy-, n-propoxy- or iso-propoxy-substituted C3-C6-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur,\nvery particularly preferably represents optionally fluorine-, chlorine-, bromine-, cyano-, nitro-, methyl-, ethyl-, n-propyl-, i-propyl-, methoxy-, ethoxy-, trifluoromethyl-, trifluoromethoxy-, methylthio-, ethylthio-, methylsulphonyl- or ethylsulphonyl-substituted phenyl,\nvery particularly preferably represents optionally fluorine-, chlorine-, bromine-, methyl-, ethyl-, n-propyl-, i-propyl-, methoxy-, ethoxy-, trifluoromethyl- or trifluoromethoxy-substituted benzyl,\nvery particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, methyl- or ethyl-substituted furanyl, thienyl or pyridyl,\nvery particularly preferably represents optionally fluorine-, chlorine-, methyl- or ethyl-substituted phenoxy-C1-C4-alkyl or\nvery particularly preferably represents respectively optionally fluorine-, chlorine-, amino-, methyl- or ethyl-substituted pyridyloxy-C1-C4-alkyl, pyrimidyloxy-C1-C4-alkyl or thiazolyloxy-C1-C4-alkyl.\nR2 very particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C14-alkyl, C2-C14-alkenyl, C1-C4-alkoxy-C2-C6-alkyl or poly-C1-C4-alkoxy-C2-C6-alkyl,\nvery particularly preferably represents optionally fluorine-, chlorine-, methyl-, ethyl-, n-propyl-, iso-propyl- or methoxy-substituted C3-C6-cycloalkyl,\nor very particularly preferably represents respectively optionally fluorine-, chlorine-, cyano-, nitro-, methyl-, ethyl-, n-propyl-, i-propyl-, methoxy-, ethoxy-, trifluoromethyl- or trifluoromethoxy-substituted phenyl or benzyl.\nR3 very particularly preferably represents optionally fluorine- or chlorine-substituted methyl, ethyl, propyl, iso-propyl, n-butyl, tert-butyl or respectively optionally fluorine-, chlorine-, bromine-, methyl-, ethyl-, iso-propyl-, tert-butyl-, methoxy-, ethoxy-, iso-propoxy-, trifluoromethyl-, trifluoromethoxy-, cyano- or nitro-substituted phenyl or benzyl.\nR4 and R5 independently of one another each very particularly preferably represent respectively optionally fluorine- or chlorine-substituted C1-C4-alkyl, C1-C4-alkoxy, C1-C4-alkylamino, di-(C1-C4-alkyl)-amino or C1-C4-alkylthio or very particularly preferably represent respectively optionally fluorine-, chlorine-, bromine-, nitro-, cyano-, methyl-, methoxy-, trifluoromethyl- or trifluoromethoxy-subsfituted phenyl, phenoxy or phenylthio.\nR6 and R7 independently of one another each very particularly preferably represent hydrogen, very particularly preferably represent respectively optionally fluorine- or chlorine-substituted C1-C4-alkyl, C3-C6-cycloalkyl, C1-C4-alkoxy, C3-C4-alkenyl or C1-C4-alkoxy-C2-C4-alkyl, very particularly preferably represent respectively optionally fluorine-, chlorine-, bromine-, methyl-, methoxy- or trifluoromethyl-substituted phenyl or benzyl, or together very particularly preferably represent an optionally methyl- or ethyl-substituted C5-C6-alkylene radical in which optionally one methylene group is replaced by oxygen or sulphur.\nParticular preference is given to compounds of the formula (I) in which Z does not represent hydrogen.\nPreference is also given to those compounds where D does not represent methyl.\nThe abovementioned definitions or illustrations of radicals mentioned generally or in preferred ranges can be combined with each other as desired, i.e. also between the respective ranges and preferred ranges. They apply correspondingly to the final products and to the precursors and intermediates.\nFor the purpose of the invention, preference is given to compounds of the formula (I) in which there exists a combination of the meanings mentioned above as preferred (preferably).\nFor the purpose of the invention, particular preference is given to compounds of the formula (I) in which there exists a combination of the meanings mentioned above as particularly preferred.\nFor the purpose of the invention, very particular preference is given to compounds of the formula (I) in which there exists a combination of the meanings mentioned above as very particularly preferred.\nSaturated or unsaturated hydrocarbon radicals such as alkyl or alkenyl can, as far as possible, in each case be straight-chain or branched, also in combination with hetero atoms, for example in alkoxy.\nOptionally substituted radicals can be mono- or polysubstituted, it being possible in the case of polysubstitution for the substituents to be identical or different.\nIn addition to the compounds mentioned in the Preparation Examples, the following compounds of the formula (I-1-a) may be specifically mentioned:\nIn addition to the compounds mentioned in the preparation examples, the following compounds of the formula (1-2-a) may be specifically mentioned:\nIn addition to the compounds mentioned in the preparation examples, the following compounds of the formula (I-3-a) may be specifically mentioned:\nIn addition to the compounds mentioned in the Preparation Examples, the following compounds of the formula (I-4-a) may be specifically mentioned:\nIf according to process (A) ethyl N-[(3-chloro-6-methyl)-phenylacetyl]-1-amino-4-ethyl-cyclohexane-carboxylate is used as starting material, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (B) ethyl O-[(2,5-dichloro)-phenylacetyl]-2-hydroxy-isobutyrate is used, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (C) ethyl 2-[(2-chloro-5-methyl)-phenyl]-4-(4-methoxy)-benzylmercapto-4-methyl-3-oxo-valerate is used, the course of the process according to the invention can be represented by the following reaction scheme: \nIf, for example, according to process (D) (chlorocarbonyl)-2-[(3-chloro-6-methyl)-phenyl]-ketene and 4-fluoro-propiophenone are used as starting materials, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (Excex1) 3-[(2,5-dichloro)-phenyl]-5,5-dimethylpyrrolidine-2,4-dione and pivaloyl chloride are used as starting materials, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (E) (variant xcex2) 3-[(2,5-dichloro)-phenyl]-4-hydroxy-5-methyl-5-phenyl-xcex943-dihydrofuran-2-one and acetic anhydride are used as starting materials, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (F) 8-[(2-chloro-5-methyl)-phenyl]-5,5-pentamethylene-pyrrolidine-2,4-dione and ethoxyethyl chloroformate are used as starting materials, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (G) 3-[(2-bromo-5-methyl)-phenyl]-4-hydroxy-5-methyl-6-(3-pyridyl)-pyrone and methyl chloromonothioformate are used as starting materials, the course of the reaction can be represented in the following manner: \nIf according to process (H) 2-[(2,5-dimethyl)-phenyl]-5,5-[(3-methyl)-pentamethylene]pyrrolidine-2,4-dione and methanesulphonyl chloride are used as starting materials, the course of the reaction can be represented by the following reaction scheme: \nIf according to process (I) 2-[(2-chlroro-5methyl)-phenyl]-4-hydroxy-5-methyl-6-(2-pyridyl)-pyrone and 2,2,2-trifluoroethyl chloromethanethio-phosphonate are used as starting materials, the course of the reaction can be represented by the following reaction scheme: \nIf according to process (J) 3-[(2,5-dichloro)-phenyl]-5-cyclopropyl-5-methyl-pyrrolidine-2,4-dione and NaOH are used as components, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (K) (variant xcex1) 3-[(2-chloro-5-methyl)-phenyl]-4-hydroxy-5,5-tetramethylene-xcex943-dihydro-furan-2-one and ethyl isocyanate are used as starting materials, the course of the reaction can be represented by the following reaction scheme: \nIf according to process (K) (variant xcex2) 3-[(2-chloro-5-methyl)-phenyl]-5,5-dimethyl-pyrrolidine-2,4-dione and dimethylcarbamoyl chloride are used as starting materials, the course of the reaction can be represented by the following scheme: \nThe compounds of the formula (II) \nin which\n5 A, B, X, Z and R8 are each as defined above,\nrequired as starting materials in process (A) according to the invention are novel.\nThe acylamino acid esters of the formula (II) are obtained, for example, when amino acid derivatives of the formula (XVIII) \nin which\nA, B and R8 are as defined above,\nare acylated using substituted phenylacetyl halides of the formula (XIX) \nxe2x80x83in which\nX and Z are each as defined above and\nHal represents chlorine or bromine,\n(Chem. Reviews 52, 237-416 (1953); Bhattacharya, Indian J. Chem. 6, 341-5, 1968)\nor when acylamino acids of the formula (XX) \nxe2x80x83in which\nA, B, X and Z are each as defined above,\nare esterified (Chem. Ind. (London) 1568 (1968)).\nThe compounds of the formula (XX) \nin which\nA, B, X and Z are each as defined above,\nare novel.\nThe compounds of the formula (XX) are obtained when amino acids of the formula (XXI) \nin which\nA and B are each as defined above,\nare acylated using substituted phenylacetyl halides of the formula (XIX) \nxe2x80x83in which\nX and Z are each as defined above and\nHal represents chlorine or bromine,\naccording to Schotten-Baumann (Organikum, VEB Deutscher Verlag der Wissenschaften, Berlin 1977, p. 505).\nSome of the compounds of the formula (XIX) are novel. They can be prepared by known methods.\nThe compounds of the formula (XIX) are obtained, for example, by reacting substituted phenylacetic acids of the formula (XXII) \nin which\nX and Z are each as defined above,\nwith halogenating agents (for example thionyl chloride, thionyl bromide, oxalyl chloride, phosgene, phosphorus trichloride, phosphorus tribromide or phosphorus pentachloride), if appropriate in the presence of a diluent (for example optionally chlorinated aliphatic or aromatic hydrocarbons such as toluene or methylene chloride) at temperatures from xe2x88x9220xc2x0 C. to 150xc2x0 C., preferably from xe2x88x9210xc2x0 C. to 100xc2x0 C.\nThe compounds of the formula (XXII) are novel with the exception of 2,5-dichlorophenylacetic acid (CAS 5398798), 5-chloro-2-methoxyphenylacetic acid (CAS 7569-6-22), 2-chloro-5-methylphenylacetic acid (CAS 81682-39-5), 2,5-difluorophenylacetic acid (CAS 85068-27-5), 2-bromo-5-methylphenylacetic acid (BRN 3 249 577) and 2-chloro-5-trifluoromethylphenylacetic acid (CAS 22893-39-6), they can be prepared by methods known from the literature (Organikum, 15th edition, p. 533, VEB Deutscher Verlag der Wissenschaften, Berlin 1977). The compounds of the formula (XMI) are obtained, for example, by hydrolysing substituted phenylacetic acid esters of the formula (XXIII) \nin which\nX, Z and R8 are each as defined above,\nat temperatures between 0xc2x0 C. and 150xc2x0 C., preferably between 20xc2x0 C. and 100xc2x0 C., in the presence of an acid (for example an inorganic acid such as hydrochloric acid) or of a base (for example of an alkali metal hydroxide such as sodium hydroxide or potassium hydroxide) and, if appropriate, of a diluent (for example of an aqueous alcohol such as methanol or ethanol).\nThe compounds of the formula (XXIII) are novel with the exception of methyl 2,5-dichlorophenylacetate (CAS 96129-66-7) and methyl 5-chloro-2-methoxy-phenylacetate (CAS 26939-01-5), they can be prepared by methods known in principle.\nThe compounds of the formula (XXIII) are obtained, for example, by reacting substituted 1,1,1-trichloro-2-phenylethanes of the formula (XXIV) \nin which\nX and Z are each as defined above,\nfirst with alkoxides (for example alkali metal alkoxides such as sodium methoxide or sodium ethoxide) in the presence of a diluent (for example the alcohol derived from the alkoxide) at temperatures between 0xc2x0 C. and 150xc2x0 C., preferably between 20xc2x0 C. and 120xc2x0 C., and then reacting with an acid (preferably an inorganic acid, such as sulphuric acid) at temperatures between xe2x88x9220xc2x0 C. and 150xc2x0 C., preferably 0xc2x0 C. and 100xc2x0 C. (cf. DE 3 314 249).\nThe compounds of the formula (XXIV) are novel, they can be prepared by methods known in principle.\nThe compounds of the formula (XXIV) are obtained, for example, when anilines of the formula (XXV) \nin which\nX and Z are as defined above,\nare reacted with vinylidene chloride (CH2=CCl2) in the presence of an alkyl nitrite of the formula (XXVI)\nR21xe2x80x94ONOxe2x80x83xe2x80x83(XXVI)\nxe2x80x83in which\nR21 represents alkyl, preferably C1-C6-alkyl,\nin the presence of copper(II) chloride and, if appropriate, in the presence of a diluent (for example of an aliphatic nitrile such as acetonitrile) at a temperature of xe2x88x9220xc2x0 C. to 80xc2x0 C., preferably 0xc2x0 C. to 60xc2x0 C.\nThe compounds of the formulae (XXV) and (XXVI) are known compounds of organic chemistry. Copper(II) chloride and vinylidene chloride have long been known and are commercially available.\nSome of the compounds of the formulae (XVIII) and (XXI) are known and/or they can be prepared by known processes (see, for example, Compagnon, Miocque Ann. Chim. (Paris) [14]5, p. 11-22, 23-27 (1970)).\nThe substituted cyclic aminocarboxylic acids of the formula (XXIa), in which A and B form a ring, are in general obtainable by the Bucherer-Bergs synthesis or by the Strecker synthesis and are in each case obtained here in different isomeric forms. Thus, according to the conditions of the Bucherer-Bergs synthesis mainly the isomers (in the following designated as xcex2 for the sake of simplicity) in which the radicals R and the carboxyl group are equatorial are obtained, while according to the conditions of the Strecker synthesis mainly the isomers (in the following designated as xcex1 for the sake of simplicity) are obtained in which the amino group and the radicals R are equatorial. \n(L. Munday, J. Chem. Soc. 4372 (1961); J. T. Eward, C. Jitrangeri, Can. J. Chem. 53, 3339 (1975).\nFurthermore the starting materials of the formula (II) \nin which\nA, B, X, Z and R8 are each as defined above,\nused in the above process (A) can be prepared by reacting aminonitriles of the formula (XXVII) \nxe2x80x83in which\nA and B are each as defined above,\nwith substituted phenylacetyl halides of the formula (XIX) \nxe2x80x83in which\nX, Z and Hal are each as defined above,\nto give compounds of the formula (XXVIII) \nxe2x80x83in which\nA, B, X and Z are each as defined above,\nand then subjecting these to an acidic alcoholysis.\nThe compounds of the formula (XXVIII) are also novel.\nThe compounds of the formula (III) \nin which\nA, B, X, Z and R8 are each as defined above,\nrequired as starting materials in process (B) according to the invention are novel.\nThey can be prepared in a simple manner by methods known in principle.\nThe compounds of the formula (III) are obtained, for example, when 2-hydroxycarboxylic acid esters of the formula (XXIX) \nin which\nA, B and R8 are each as defined above,\nare acylated using substituted phenylacetyl halides of the formula (XIX) \nxe2x80x83in which\nX, Z and Hal are each as defined above\n(Chem. Reviews 52, 237-416 (1953)).\nThe compounds of the formula (IV) \nin which\nA, B, W, X, Z and R8 are each as defined above,\nrequired as starting materials in the above process (C) are novel.\nThey can be prepared by methods known in principle.\nThe compounds of the formula (IV) are obtained, for example, when substituted phenylacetic acid esters of the formula (XXIII) \nin which\nX, R8 and Z are each as defined above,\nare acylated using 2-benzylthio-carbonyl halides of the formula (XXX) \nin which\nA, B and W are each as defined above and\nHal represents halogen (in particular chlorine or bromine),\nin the presence of strong bases (see, for example, M. S. Chambers, E. J. Thomas, D. J. Williams, J. Chem. Soc. Chem. Commun., (1987), 1228).\nThe benzylthio-carbonyl halides of the formula (XXX) are known in some cases and/or can be prepared by known methods (J. Antibiotics (1983), 26, 1589).\nThe halogenocarbonylketenes of the formula (VI) in which Z does not represent hydrogen, which are required as starting materials in process (D), are novel. They can be prepared in a simple manner by methods known in principle (cf., for example, Org. Prep. Proced. Int., 7, (4), 155-158, 1975 and DE 1 945 703).\nThe compounds of the formula (VI) \nin which\nX and Z are each as defined above and\nHal represents chlorine or bromine,\nare obtained when substituted phenylmalonic acids of the formula (XXXI) \nxe2x80x83in which\nX and Z are each as defined above,\nare reacted with acid halides, for example thionyl chloride, phosphorus(V) chloride, phosphorus(III) chloride, oxalyl chloride, phosgene or thionyl bromide, if appropriate in the presence of catalysts, for example diethylformamide, methylstearylformamide or triphenylphosphine and, if appropriate, in the presence of bases, for example pyridine or triethylamine, at a temperature between xe2x88x9220xc2x0 C. and 200xc2x0 C., preferably between 0xc2x0 C. and 150xc2x0 C.\nThe substituted phenylmalonic acids of the formula (XXXI) in which Z does not represent hydrogen are novel. However, they may be prepared in a simple manner by known processes (cf., for example, Organikum, VEB Deutscher Verlag der Wissenschaften, Berlin 1977, p. 517 ff), for example by hydrolysis of substituted phenylmalonic esters of the formula (XXXI) \nin which\nX, Z and R8 are each as defined above.\nThe carbonyl compounds of the formula (V) or their silyl enole ethers of the formula (Va) \nin which\nA, D and R8 xe2x80x2are each as defined above,\nrequired as starting materials for process (D) according to the invention are compounds which are commercially available, generally known or accessible by known processes.\nThe malonic acid esters of the formula (XXXII) \nin which\nR8, X and Z are each as defined above, and Z is not hydrogen, are novel.\nThey can be prepared by generally known methods of organic chemistry (cf., for example, Tetrahedron Lett. 27, 2763 (1986) and Organikum VEB Deutscher Verlag der Wissenschaften, Berlin 1977, p.587 ff).\nThe acid halides of the formula (VII), carboxylic anhydrides of the formula (VIII), chloroformic acid esters or chloroformic acid thioesters of the formula (IX), chloromonothioformic acid esters or chlorodithioformic acid esters of the formula (X), sulphonyl chlorides of the formula (XII), phosphorus compounds of the formula (XIII) and metal hydroxides, metal alkoxides or amines of the formula (XIV) and (XV) and isocyanates of the formula (XVI) and carbamoyl chlorides of the formula (XVII) additionally required as starting materials for carrying out processes (F), (G), (H), (I), (J) and (K) according to the invention are generally known compounds of organic or inorganic chemistry.\nThe compounds of the formulae (V), (VII) to (XVII), (XVIII), (XXI), (XXII), (XXIX), (XXX) and (XXXI) are moreover disclosed in the patent applications cited at the outset and/or can be prepared by the methods given there.\nProcess (A) is characterized in that compounds of the formula (II) in which A, B, X, Z and R8 are each as defined above are subjected to an intramolecular condensation in the presence of a diluent and in the presence of a base.\nSuitable diluents for the process (A) according to the invention are all organic solvents which are inert to the reaction participants. Those preferably utilizable are hydrocarbons, such as toluene and xylene, furthermore ethers, such as dibutyl ether, tetrahydrofuran, dioxane, glycol dimethyl ether and diglycol dimethyl ether, additionally polar solvents, such as dimethyl sulphoxide, sulpholane, dimethyl-formamide and n-methyl-pyrrolidone, and also alcohols such as methanol, ethanol, propanol, iso-propanol, butanol, iso-butanol and tert-butanol.\nSuitable bases (deprotonating agents) for carrying out process (A) according to the invention are all customary proton acceptors. Those preferably utilizable are alkali metal and alkaline earth metal oxides, hydroxides and carbonates, such as sodium hydroxide, potassium hydroxide, magnesium oxide, calcium oxide, sodium carbonate, potassium carbonate and calcium carbonate, each of which can also be employed in the presence of a phase-transfer catalyst, for example triethylbenzylammonium chloride, tetrabutylammonium bromide, Adogene 464 (=methyltrialkyl(C8-C10)ammonium chloride) or TDA 1 (=tris-(methoxyethoxyethyl)-amine. Alkali metals such as sodium or potassium can also be used. Furthermore, alkali metal and alkaline earth metal amides and hydrides, such as sodium amide, sodium hydride and calcium hydride, and additionally also alkali metal alkoxides, such as sodium methoxide, sodium ethoxide and potassium tert-butoxide can be employed.\nWhen carrying out process (A) according to the invention, the reaction temperature can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between 0xc2x0 C. and 250xc2x0 C., preferably between 50xc2x0 C. and 150xc2x0 C.\nProcess (A) according to the invention is generally carried out under atmospheric pressure.\nWhen carrying out process (A) according to the invention, the reaction component of the formula (II) and the deprotonating base are generally employed in equimolar to approximately double equimolar amounts. However, it is also possible to use one component or the other in a relatively large excess (up to 3 mol).\nProcess (B) is characterized in that compounds of the formula (III) in which A, B, X, Z and R8 are each as defined above are condensed intramolecularly in the presence of a diluent and in the presence of a base.\nSuitable diluents for the process (B) according to the invention are all organic solvents which are inert to the reaction participants. Those preferably utilizable are hydrocarbons, such as toluene and xylene, furthermore ethers, such as dibutyl ether, tetrahydrofuran, dioxane, glycol dimethyl ether and diglycol dimethyl ether, and additionally polar solvents, such as dimethyl sulphoxide, sulpholane, dimethylformamide and N-methyl-pyrrolidone. Alcohols such as methanol, ethanol, propanol, iso-propanol, butanol, iso-butanol and tert-butanol can also be used.\nSuitable bases (deprotonating agents) for carrying out the process (B) according to the invention are all customary proton acceptors. Those preferably utilizable are alkali metal and alkaline earth metal oxides, hydroxides and carbonates, such as sodium hydroxide, potassium hydroxide, magnesium oxide, calcium oxide, sodium carbonate, potassium carbonate and calcium carbonate, each of which can also be employed in the presence of phase-transfer catalysts, for example triethylbenzylammonium chloride, tetrabutylammonium bromide, Adogene 464 (=methyltrialkyl(C8-C10-ammonium chloride) or TDA 1 (=tris-(methoxyethoxyethyl)-amine). Alkali metals such as sodium or potassium can also be used. Suitable are also alkali metal and alkaline earth metal amides and hydrides, such as sodium amide, sodium hydride and calcium hydride, and also alkali metal alkoxides, such as sodium methoxide, sodium ethoxide and potassium tert-butoxide.\nWhen carrying out process (B) according to the invention, the reaction temperature can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between 0xc2x0 C. and 250xc2x0 C., preferably between 50xc2x0 C. and 150xc2x0 C.\nProcess (B) according to the invention is generally carried out under atmospheric pressure.\nWhen carrying out the process (B) according to the invention, the reaction components of the formula (III) and the deprotonating bases are generally employed in approximately equimolar amounts. However, it is also possible to use one component or the other in a relatively large excess (up to 3 mol).\nProcess (C) is characterized in that compounds of the formula (IV) in which A, B, W, X, Z and R8 are each as defined above are cyclized intramolecularly in the presence of an acid and, if appropriate, in the presence of a diluent.\nSuitable diluents for the process (C) according to the invention are all organic solvents which are inert to the reaction participants. Those preferably utilizable are hydrocarbons, such as toluene and xylene, furthermore halogenated hydrocarbons, such as dichloromethane, chloroform, ethylene chloride, chlorobenzene, dichlorobenzene, and additionally polar solvents, such as dimethyl sulphoxide, sulpholane, dimethylformamide and N-methyl-pyrrolidone. Alcohols such as methanol, ethanol, propanol, iso-propanol, butanol, iso-butanol and tert-butanol can also be used.\nThe acid employed can, if appropriate, also be used as a diluent.\nAcids which can be employed in process (C) according to the invention are all customary inorganic and organic acids, for example hydrohalic acids, sulphuric acid, alkyl-, aryl- and haloalkyl sulphonic acids; halogenated alkylcarboxylic acids, for example trifluoroacetic acid, are used in particular.\nWhen carrying out process (C) according to the invention, the reaction temperature can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between 0xc2x0 C. and 250xc2x0 C., preferably between 50xc2x0 C. and 150xc2x0 C.\nProcess (C) according to the invention is generally carried out under atmospheric pressure.\nWhen carrying out process (C) according to the invention, the reaction components of the formulae (IV) and the acid are employed, for example, in equimolar amounts. However, it is, if appropriate, also possible to employ the acid in catalytic amounts.\nProcess (D) according to the invention is characterized in that carbonyl compounds of the formula (V) or their silyl enol ethers of the formula (Va) are reacted with ketene acid halides of the formula (VI), if appropriate in the presence of a diluent and if appropriate in the presence of an acid acceptor.\nSuitable diluents for the process (D) according to the invention are all organic solvents which are inert to the reaction participants. Those preferably utilizable are hydrocarbons, such as o-dichlorobenzene, tetraline, toluene and xylene, furthermore ethers, such as dibutyl ether, glycol dimethyl ether and diglycol dimethyl ether, and additionally polar solvents, such as dimethyl sulphoxide, sulpholane, dimethylformamide or N-methyl-pyrrolidone.\nAcid acceptors which can be used when carrying out process (D) according to the invention are all customary acid acceptors.\nThose preferably utilizable are tertiary amines, such as triethylamine, pyridine, diazabicyclooctane (DABCO), diazabicycloundecene (DBU), diazabicyclononene (DBN), Hunig base or N,N-dimethyl-aniline.\nWhen carrying out process (D) according to the invention, the reaction temperature can be varied within a relatively wide range. The reaction is expediently carried out at temperatures between 0xc2x0 C. and 250xc2x0 C., preferably between 50xc2x0 C. and 220xc2x0 C.\nProcess (D) according to the invention is preferably carried out under atmospheric pressure.\nWhen carrying out process (D) according to the invention, the reaction components of the formulae (V) and (VI) and, if appropriate, the acid acceptor are in general employed in approximately equimolar amounts. However, it is also possbile to use one component or the other in a relatively large excess (up to 5 mol).\nProcess (Excex1) is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with carboxylic acid halides of the formula (VII), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nSuitable diluents for the process (Eac) according to the invention are all solvents inert to the acid halides. Those preferably utilizable are hydrocarbons, such as benzine, benzene, toluene, xylene and tetraline, furthermore halogenated hydrocarbons, such as methylene chloride, chloroform, carbon tetrachloride, chlorobenzene and o-dichlorobenzene, and also ketones, such as acetone and methyl isopropyl ketone, furthermore ethers, such as diethyl ether, tetrahydrofuran and dioxan, moreover carboxylic acid esters, such as ethyl acetate, and also strongly polar solvents, such as dimethylformamide, dimethylsulphoxide and sulpholane. If the stability to hydrolysis of the acid halide permits, the reaction can also be carried out in the presence of water.\nSuitable acid-binding agents in the reaction of process (Es) according to the invention are all customary acid acceptors. Those preferably utilizable are tertiary amines, such as triethylamine, pyridine, diazabicyclooctane (DABCO), diazabicycloundecene (DBU), diazabicyclononene (DBN), Huinig base and N,N-dimethyl-aniline, furthermore alkaline earth metal oxides, such as magnesium oxide and calcium oxide, and also alkali metal and alkaline earth metal carbonates, such as sodium carbonate, potassium carbonate and calcium carbonate and also alkali metal hydroxides such as sodium hydroxide and potassium hydroxide.\nThe reaction temperature in the process (Excex1) according to the invention can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between xe2x88x9220xc2x0 C. and +150xc2x0 C., preferably between 0xc2x0 C. and 110xc2x0 C.\nWhen carrying out process (Excex1) according to the invention, the starting materials of the formulae (I-1-a) to (I-4-a) and the carboxylic acid halide of the formula (VII) are in general each used in approximately equivalent amounts. However, it is also possible to employ the carboxylic acid halide in a relatively large excess (up to 5 mol). Work-up is carried out according to customary methods.\nProcess (Excex2) is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are each reacted with carboxylic anhydrides of the formula (VIII), if appropriate in the presence of a diluent and if appopriate in the presence of an acid-binding agent.\nPreferred diluents for the process (Excex2) according to the invention are those diluents which are also preferred when using acid halides. Otherwise, a carboxylic anhydride employed in excess may also simultaneously function as diluent.\nPossible acid-binding agents added in process (Excex2) are preferably those acid-binding agents that are also preferred when using acid-halides.\nThe reaction temperature in the process (Excex2) according to the invention can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between xe2x88x9220xc2x0 C. and +150xc2x0 C., preferably between 0xc2x0 C. and 100xc2x0 C.\nWhen carrying out process (Excex2) according to the invention, the starting materials of the formulae (I-1-a) to (I-4-a) and the carboxylic anhydride of the formula (VIII) are in general each used in approximately equivalent amounts. However, it is also possible to employ the carboxylic anhydride in a relatively large excess (up to 5 mol). Work-up is carried out according to customary methods.\nIn general, a procedure is used in which diluent and excess carboxylic anhydride and the resulting carboxylic acid are removed by distillation or by washing with an organic solvent or with water.\nProcess (F) is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with chloroformic acid esters or chloroformic acid thioesters of the formula (IX), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nSuitable acid-binding agents for process (F) according to the invention are all customary acid acceptors. Those preferably utilizable are tertiary amines, such as triethylamine, pyridine, DABCO, DBU, DBA, Hunig base and N,N-dimethyl-aniline, furthermore alkaline earth metal oxides, such as magnesium oxide and calcium oxide, additionally alkali metal and alkaline earth metal carbonates, such as sodium carbonate, potassium carbonate and calcium carbonate and also alkali metal hydroxides such as sodium hydroxide and potassium hydroxide.\nSuitable diluents for the process (F) according to the invention are all solvents which are inert to the chloroformic acid esters or chloroformic acid thioesters. Those preferably utilizable are hydrocarbons, such as benzine, benzene, toluene, xylene and tetraline, furthermore halogenated hydrocarbons, such as methylene chloride, chloroform, carbon tetrachloride, chlorobenzene and o-dichlorobenzene, additionally ketones, such as acetone and methyl isopropyl ketone, furthermore ethers, such as diethyl ether, tetrahydrofuran and dioxan, moreover carboxylic acid esters, such as ethyl acetate, furthermore nitriles, such as acetonitrile, and also strongly polar solvents, such as dimethylformamide, dimethyl sulphoxide and sulpholane.\nWhen carrying out process (F) according to the invention, the reaction temperature can be varied within a relatively wide range. The reaction temperature is generally between xe2x88x9220xc2x0 C. and +100xc2x0 C., preferably between 0xc2x0 C. and 50xc2x0 C.\nProcess (F) according to the invention is in general carried out under atmospheric pressure.\nWhen carrying out process (F) according to the invention, the starting materials of the formulae (I-1-a) to (I-4-a) and the appropriate chloroformic acid ester or chloroformic acid thioester of the formula (IX) are in general each used in approximately equivalent amounts. However, it is also possible to employ one component or the other in a relatively large excess (up to 2 mol). Work-up is carried out according to customary methods. In general, a procedure is used in which precipitated salts are removed and the reaction mixture which remains is concentrated by stripping off the diluent.\nProcess (G) according to the invention is characterized in that compounds of the formula (I-1-a) to (I-4-a) are in each case reacted with compounds of the formula (X) in the presence of a diluent and, if appropriate, in the presence of an acid-binding agent.\nIn preparation process (G), about 1 mol of chloromonothioformic acid ester or chlorodithioformic acid ester of the formula (X) is reacted per mol of a starting material of the formulae (I-1-a) to (I-4-a), at 0 to 120xc2x0 C., preferably at 20 to 60xc2x0 C.\nDiluents which may be added, if appropriate, are all inert polar organic solvents, such as ethers, amides, sulphones, sulphoxides, and also halogenoalkanes.\nDimethyl sulphoxide, tetrahydrofuran, dimethylformamide or methylene chloride are preferably employed.\nIf, in a preferred embodiment, the enolate salt of the compounds (I-1-a) to (I-4-a) is prepared by addition of strong deprotonating agents, for example sodium hydride or potassium tert-butoxide, the further addition of acid-binding agents can be dispensed with.\nIf acid-binding agents are used, customary inorganic or organic bases are suitable; sodium hydroxide, sodium carbonate, potassium carbonate, pyridine and triethylamine may be mentioned by way of example.\nThe reaction can be carried out at atmospheric pressure or at elevated pressure; it is preferably carried out at atmospheric pressure. Work-up takes place according to customary methods.\nProcess (H) according to the invention is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with sulphonyl chlorides of the formula (XII), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nIn preparation process (H), about 1 mol of sulphonyl chloride of the formula (XII) is reacted per mol of starting material of the formulae (I-1-a) to (I-4-a) at xe2x88x9220 to 150xc2x0 C., preferably at 0 to 70xc2x0 C.\nProcess (H) is preferably carried out in the presence of a diluent.\nSuitable diluents are all inert polar organic solvents, such as ethers, amides, ketones, carboxylic acid esters, nitriles, sulphones, sulphoxides or halogenated hydrocarbons such as methylene chloride.\nDimethyl sulphoxide, tetrahydrofuran, dimethylformamide and methylene chloride are preferably employed.\nIf, in a preferred embodiment, the enolate salt of the compounds (I-1-a) to (I-4-a) is prepared by addition of strong deprotonating agents (for example sodium hydride or potassium tert-butoxide), the further addition of acid-binding agents can be dispensed with.\nIf acid-binding agents are employed, customary inorganic or organic bases are suitable; sodium hydroxide, sodium carbonate, potassium carbonate, pyridine and triethylamine may be mentioned by way of example.\nThe reaction can be carried out at atmospheric pressure or at elevated pressure; it is preferably carried out at atmospheric pressure. Work-up takes place according to customary methods.\nThe process (I) according to the invention is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with phosphorus compounds of the formula (XIII), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nIn preparation process (I), 1 to 2, preferably 1 to 1.3, mol of the phosphorus compound of the formula (XIII) is reacted per 1 mol of the compounds (I-1-a) to (I-4-a) at temperatures between xe2x88x9240xc2x0 C. and 150xc2x0 C., preferably between xe2x88x9210xc2x0 C. and 110xc2x0 C., to give compounds of the formulae (I-1-e) to (I-4-e).\nThe process (I) is preferably carried out in the presence of a diluent.\nSuitable diluents are all inert polar organic solvents, such as ethers, carboxylic acid esters, halogenated hydrocarbons, ketones, amides, nitriles, sulphones, sulphoxides, etc.\nAcetonitrile, dimethyl sulphoxide, tetrahydrofuran, dimethylformamide or methylene chloride are preferably employed.\nAcid-binding agents which may be added, if appropriate, are customary inorganic or organic bases, such as hydroxides, carbonates or amines. By way of example, sodium hydroxide, sodium carbonate, potassium carbonate, pyridine and triethylamine may be mentioned.\nThe reaction may be carried out at atmospheric pressure or at elevated pressure; it is preferably carried out at atmospheric pressure. Work-up takes place according to conventional methods of organic chemistry. The end products are preferably purified by crystallization, chromatographic purification or by so-called xe2x80x9cincipient distillationxe2x80x9d, i.e. removal of the volatile constituents in vacuo.\nProcess (J) is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with metal hydroxides or metal alkoxides of the formula (XIV) or amines of the formula (XV), if appropriate in the presence of a diluent.\nPreferred diluents for process (J) according to the invention are ethers such as tetrahydrofuran, dioxan and diethyl ether or else alcohols such as methanol, ethanol and isopropanol, but also water. Process (J) according to the invention is generally carried out under atmospheric pressure. The reaction temperature is in general between xe2x88x9220xc2x0 C. and 100xc2x0 C., preferably between 0xc2x0 C. and 50xc2x0 C.\nProcess (K) according to the invention is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with (Kxcex1) compounds of the formula (XVI), if appropriate in the presence of a diluent and if appropriate in the presence of a catalyst, or (Kxcex2) with compounds of the formula (XVII), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nIn preparation process (Kxcex1), about 1 mol of isocyanate of the formula (XVI) is reacted per mole of starting material of the formulae (I-1-a) to (I-4-a) at 0 to 100xc2x0 C., preferably at 20 to 50xc2x0 C.\nProcess (Kxcex1) is preferably carried out in the presence of a diluent.\nSuitable diluents are all inert organic solvents, such as aromatic hydrocarbons, halogenated hydrocarbons, ethers, amides, nitriles, sulphones or sulphoxides.\nCatalysts may, if desired, be added to accelerate the reaction. The catalysts employed can very advantageously be organotin compounds, for example dibutyltin dilaurate.\nThe reaction is preferably carried out at atmospheric pressure.\nIn preparation process (Kxcex2), about 1 mol of carbamoyl chloride of the formula (XVII) is reacted at 0 to 150xc2x0 C., preferably at 20 to 70xc2x0 C., per mole of starting material of the formulae (I-1-a) to (I-4-a).\nPossible diluents optionally added are all inert polar organic solvents, such as ethers, carboxylic acid esters, nitriles, ketones, amides, sulphones, sulphoxides or halogenated hydrocarbons.\nDimethyl sulphoxide, tetrahydrofuran, dimethylformamide or methylene chloride are preferably employed.\nIf, in a preferred embodiment, the enolate salt of the compound (I-1-a) to (I-4-a) is prepared by addition of strong deprotonating agents (e.g. sodium hydride or potassium tertiary butoxide), the further addition of acid-binding agents can be dispensed with.\nIf acid-binding agents are employed, customary inorganic or organic bases are suitable; those which may be mentioned by way of example are sodium hydroxide, sodium carbonate, potassium carbonate, triethylamine or pyridine.\nThe reaction can be carried out at atmospheric pressure or at elevated pressure, preferably at atmospheric pressure. Work-up takes place according to customary methods.\nThe active compounds are suitable for controlling animal pests, preferably arthropods and nematodes, in particular insects and arachnida, which are encountered in agriculture, in forestry, in the protection of stored products and of materials, and in the hygiene field. They are active against normally sensitive and resistant species and against all or some stages of development. The abovementioned pests include:\nFrom the order of Isopoda, for example, Oniscus asellus, Armadillidium vulgare and Porcellio scaber. \nFrom the order of the Diplopoda, for example, Blaniulus guttulatus. \nFrom the order of the Chilopoda, for example, Geophilus carpophagus and Scutigera spec.\nFrom the order of the Symphyla, for example, Scutigerella immaculata. \nFrom the order of the Thysanura, for example, Lepisma saccharina. \nFrom the order of the Collembola, for example, Onychiurus armatus. \nFrom the order of the Orthoptera, for example, Blatta orientalis, Periplaneta americana, Leucophaea maderae, Blatella germanica, Acheta domesticus, Gryllotalpa spp., Locusta migratoria migratorioides, Melanoplus differentialis and Schistocerca gregaria. \nFrom the order of the Dermaptera, for example, Forficula auricularia. \nFrom the order of the Isoptera, for example, Reticulitermes spp.\nFrom the order of the Anoplura, for example, Phylloxera vastatrix, Pemphigus spp., Pediculus humanus corporis, Haematopinus spp. and Linognathus spp.\nFrom the order of the Mallophaga, for example, Trichodectes spp. and Damalinea spp.\nFrom the order of the Thysanoptera, for example, Hercinothrips femoralis and Thrips tabaci. \nFrom the order of the Heteroptera, for example, Eurygaster spp., Dysdercus intermedius, Piesma quadrata, Cimex lectularius, Rhodnius prolixus and Triatoma spp.\nFrom the order of the Homoptera, for example, Aleurodes brassicae, Bemisia tabaci, Trialeurodes vaporariorum, Aphis gossypii, Brevicoryne brassicae, Cryptomyzus ribis, Aphis fabae, Doralis pomi, Eriosoma lanigerum, Hyalopterus arundinis, Macrosiphum avenae, Myzus spp., Phorodon humuli, Rhopalosiphum padi, Empoasca spp., Euscelis bilobatus, Nephotettix cincticeps, Lecanium corni, Saissetia oleae, Laodelphax striatellus, Nilaparvata lugens, Aonidiella aurantii, Aspidiotus hederae, Pseudococcus spp. and Psylla spp.\nFrom the order of the Lepidoptera, for example, Pectinophora gossypiella, Bupalus piniarius, Cheimatobia brumata, Lithocolletis blancardella, Hyponomeuta padella, Plutella maculipennis, Malacosoma neustria, Euproctis chrysorrhoea, Lymantria spp, Bucculatrix thurberiella, Phyllocnistis citrella, Agrotis spp., Euxoa spp., Feltia spp., Earias insulana, Heliothis spp., Spodoptera exigua, Mamestra brassicae, Panolis flammea, Prodenia litura, Spodoptera spp., Trichoplusia ni, Carpocapsa pomonella, Pieris spp., Chilo spp., Pyrausta nubilalis, Ephestia kuehniella, Galleria mellonella, Tineola bisselliella, Tinea pellionella, Hofmannophila pseudospretella, Cacoecia podana, Capua reticulana, Choristoneura fumiferana, Clysia ambiguella, Homona magnanima and Tortrix viridana. \nFrom the order of the Coleoptera, for example, Anobium punctatum, Rhizopertha dominica, Acanthoscelides obtectus, Hylotrupes bajulus, Agelastica alni, Leptinotarsa decemlineata, Phaedon cochleariae, Diabrotica spp., Psylliodes chrysocephala, Epilachna varivestis, Atomaria spp., Oryzaephilus surinamensis, Anthonomus spp., Sitophilus spp., Otiorrhynchus sulcatus, Cosmopolites sordidus, Ceuthorrhynchus assimilis, Hypera postica, Dermestes spp., Trogoderma spp., Anthrenus spp., Attagenus spp., Lyctus spp., Meligethes aeneus, Ptinus spp., Niptus hololeucus, Gibbium psylloides, Tribolium spp., Tenebrio molitor, Agriotes spp., Conoderus spp., Melolontha melolontha, Amphimallon solstitialis and Costelytra zealandica. \nFrom the order of the Hymenoptera, for example, Diprion spp., Hoplocampa spp., Lasius spp., Monomorium pharaonis and Vespa spp.\nFrom the order of the Diptera, for example, Aedes spp., Anopheles spp., Culex spp., Drosophila melanogaster, Musca spp., Fannia spp., Calliphora erythrocephala, Lucilia spp., Chrysomyia spp., Cuterebra spp., Gastrophilus spp., Hyppobosca spp., Stomoxys spp., Oestrus spp., Hypoderma spp., Tabanus spp., Tannia spp., Bibio hortulanus, Oscinella frit, Phorbia spp., Pegomyia hyoscyami, Ceratitis capitata, Dacus oleae and Tipula paludosa. \nFrom the order of the Siphonaptera, for example, Xenopsylla cheopis and Ceratophyllus spp.\nFrom the order of the Arachnida, for example, Scorpio maurus and Latrodectus mactans. \nFrom the order of the Acarina, for example, Acarus siro, Argas spp., Ornithodoros spp., Dermanyssus gallinae, Eriophyes ribis, Phyllocoptruta oleivora, Boophilus spp., Rhipicephalus spp., Amblyomma spp., Hyalomma spp., Ixodes spp., Psoroptes spp., Chorioptes spp., Sarcoptes spp., Tarsonemus spp., Bryobia praetiosa, Panonychus spp. and Tetranychus spp.\nThe active compounds according to the invention are distinguished by a high insecticidal and acaricidal activity.\nThey can be used to particularly good effect for controlling insects which are injurious to plants, such as, for example, against the larvae of the mustard beetle (Phaedon cochleariae), against the larvae of the green rice leaf hopper (Nephotettix cincticeps) or against the caterpillars of the cabbage moth (Plutella maculipennis) (cf. the Use Examples).\nThe active compounds can be converted into the customary formulations, such as solutions, emulsions, wettable powders, suspensions, powders, dusting agents, pastes, soluble powders, granules, suspension-emulsion concentrates, natural and synthetic materials impregnated with active compound, and very fine capsules in polymeric substances.\nThese formulations are produced in a known manner, for example by mixing the active compounds with extenders, that is liquid solvents and/or solid carriers, optionally with the use of surface-active agents, that is emulsifying agents and/or dispersing agents and/or foam-forming agents.\nIn the case of the use of water as an extender, organic solvents can, for example, also be used as auxiliary solvents. As liquid solvents, there are suitable in the main: aromatics, such as xylene, toluene or alkylnaphthalenes, chlorinated aromatics and chlorinated aliphatic hydrocarbons, such as chlorobenzenes, chloroethylenes of methylene chloride, aliphatic hydrocarbons, such as cyclohexane or paraffins, for example petroleum fractions, mineral and vegetable oils, alcohols, such as butanol or glycol as well as their ethers and esters, ketones, such as acetone, methyl ethyl ketone, methyl isobutyl ketone or cyclohexanone, strongly polar solvents, such as dimethylformamide and dimethyl sulphoxide, as well as water.\nAs solid carriers there are suitable:\nfor example ammonium salts and ground natural minerals, such as kaolins, clays, talc, chalk, quartz, attapulgite, montmorillonite or diatomaceous earth, and ground synthetic minerals, such as highly disperse silica, alumina and silicates, as solid carriers for granules there are suitable: for example crushed and fractionated natural rocks such as calcite, marble, pumice, sepiolite and dolomite as well as synthetic granules of inorganic and organic meals, and granules of organic material such as sawdust, coconut shells, maize cobs and tobacco stalks; as emulsifying and/or foam-forming agents there are suitable: for example non-ionic and anionic emulsifiers, such as polyoxyethylene fatty acid esters, polyoxyethylene fatty alcohol ethers, for example alkylaryl polyglycol ethers, alkylsulphonates, alkylsulphates, arylsulphonates as well as albumen hydrolysis products, as dispersing agents there are suitable: for example lignin-sulphite waste liquors and methylcellulose.\nAdhesives such as carboxymethylcellulose and natural and synthetic polymers in the form of powders, granules or latices, such as gum arabic, polyvinyl alcohol and polyvinyl acetate, as well as natural phospholipids such as cephalins and lecithins, and synthetic phospholipids, can be used in the formulations. Further additives can be mineral and vegetable oils.\nIt is possible to use colorants such as inorganic pigments, for example iron oxide, titanium oxide and Prussian Blue, and organic dyestuffs, such as alizarin dyestuffs, azo dyestuffs and metal phthalocyanine dyestuffs, and trace nutrients such as salts of iron, manganese, boron, copper, cobalt, molybdenum and zinc.\nThe formulations in general contain between 0.1 and 95 per cent by weight of active compound, preferably between 0.5 and 90%.\nThe active compound according to the invention can be present in its commercially available formulations and in the use forms prepared from these formulations, as a mixture with other active compounds, such as insecticides, attractants, sterilizing agents, acaricides, nematicides, fungicides, growth-regulating substances or herbicides. The insecticides include, for example, phosphates, carbamates, carboxylates, chlorinated hydrocarbons, phenylureas and substances produced by microorganisms."}
-{"text": "The present invention relates to an incremental encoder and, more particularly, to an incremental encoder having a rotor provided with a plurality of indices and capable of detecting a reference or zero point without turning the rotor one full turn.\nEncoders have been widely used for electrically measuring distances and angles. Optical encoders and magnetic encoders are quite analogous in basic configuration.\nRotary encoders have been prevalently used for the electrical measurement of angles. In particular, optical encoders incorporate advanced optical techniques, they are capable of measurement with high accuracy and high resolution, they have high resolving power and are immune to disturbances such as external magnetism, and they have a long lifetime owing to their noncontact measuring operation.\nOptical encoders having such exemplary features are used in surveying instruments for measuring angles, for example.\nOptical encoders employed in current surveying instruments are classified into those of absolute type, i.e., absolute optical encoders, and those of incremental type, i.e., incremental optical encoders.\nIn an absolute optical encoder, angle and position on a circumference are in one-to-one correspondence, and positions on a circumference are registered as absolute addresses. Therefore, position information about any position can be obtained. However, an absolute optical encoder has a complex construction, and it is very difficult to construct an absolute optical encoder in a compact, lightweight manner for a surveying instrument.\nAs shown in FIG. 6, an incremental optical encoder comprises a rotor 9100 provided with a main scale 9110, a stator 9200 provided with a subscale 9210, and a detecting means having components 9310, 9320 and 9330. The main scale 9110 formed on the rotor 9100 has graduation lines arranged at equal angular pitch, e.g. of 80 sec, at the periphery of the rotor 9100 in an occulting pattern. The subscale 9210 formed on the stator 9200 has graduation lines arranged at an angular pitch equal to that of the main scale 9110 in an occulting pattern.\nThe detecting means comprises an LED (light emitting diode) 9310, a collimator lens 9320 and a photosensor 9330. The LED 9310 and the photosensor 9330 are disposed with the rotor 9100 and the stator 9200 therebetween. When the rotor 9100 rotates, the light emitted by the LED 9310 is interrupted once every time the rotor 9100 turns through an angle corresponding to one graduation of the main scale 9110, and the photosensor 9330 provides an electric signal indicating the interruption of the light. The electric signals provided by the photosensor 9330 are counted to detect the angle through which the rotor 9100 has turned.\nThe incremental optical encoder is able to start counting the output signals of the photosensor 9330 from any position of the rotor 9100 and to measure the angle of turn of the rotor 9100 from the position where counting the output signals of the photosensor 9330 is started.\nA surveying instrument employing an incremental optical encoder will be described with reference to FIGS. 7(a) and 7(b). A surveying instrument 10000 comprises a base unit 8100, a standard unit 8200 rotatably mounted on the base unit 8100 in a horizontal plane, a sighting telescope 8300 supported on the standard unit 8200 for turning in a vertical plane, a first detector 8400 for detecting the horizontal angle of the standard unit 8200, and a second detector 8500 for detecting the elevation angle of the sight line of the telescope 8300.\nThe base unit 8100 is connected by leveling screws 8160 to a leveling plate 8150 to be fixedly mounted on a tripod or the like. The level of the surveying instrument 10000 can be adjusted by turning the leveling screws 8160. The base unit 8100 is provided with a lower adjusting knob 8120 and a lower clamp knob 8130 to adjust and fix the base unit 8100. The standard unit 8200 is provided with an upper adjusting knob 8220 and an upper clamp knob 8230 to adjust and fix the standard unit 8200. The sighting telescope 8300 is provided with an elevation adjusting knob 8320 and an elevation clamp knob 8330 to adjust the elevation angle in the sighting direction of the sighting telescope 8300 and to fix the sighting telescope 8300 in an adjusted sighting direction.\nAn optical encoder included in the second detector 8500 detects the elevation angle of the sighting direction of the sighting telescope 8300 with respect to zenith, for example. The optical encoder of the second detector 8500 is provided with an index for zero detection to determine an angle from a reference. The signals are counted using the index as a reference or zero point to measure the angle. An optical encoder included in the first detector 8400 for measuring the horizontal angle of the standard unit 8200 has no particular reference direction and hence does not need any reference point. Therefore, the optical encoder is an ordinary one not provided with any index.\nThe incremental optical encoder provided with an index of the second detector 8500 will be described with reference to FIG. 8. The incremental optical encoder of the second detector 8500 comprises a rotor 8510, a stator 8520, and an optical detector comprising components 8531-8533 and 8535-8537 disposed with the rotor 8510 and the stator 8520 there between as shown. The rotor 8510 is provided at its periphery with a main scale 8511 having graduation lines formed at equal angular pitch, and a zero detection index 8512. The stator 8520 is provided with a first subscale 8521 for use in combination with the main scale 8511, and a second subscale 8522 for use in combination with the zero detecting index 8512. The optical detector 8530 comprises an index detecting unit and a main scale detecting unit. The index detecting unit comprises a first light emitting device 8531, a first collimator lens 8532 and a first photosensor 8533 and is capable of detecting the zero detecting index 8512 of the rotor 8510. The main scale detecting unit comprises a second light emitting device 8536, a second collimator lens 8535 and a second photosensor 8537. The main scale detecting unit detects light pulses produced by the occulting pattern of the main scale 8511 of the rotor 8510, and the second photosensor 8537 converts the light pulses into corresponding electric signals. The electric signals are counted to determine an angle from a zero detection point.\nA method of using the conventional surveying instrument 10000 thus constructed will be described. The leveling plate 8150 of the surveying instrument 10000 is mounted on a tripod, and the leveling plate 8150 is leveled by turning the leveling screws 8160. Then, a main switch is closed and the sighting telescope 8300 is turned one full turn to complete preparations. When the sighting telescope 8300 is turned, the rotor 8510 of the second detector 8500 is turned, and the zero detecting index 8512 sets a zero point to measure an angle from the zero point.\nThis conventional surveying instrument 10000 needs to turn by one full turn the sighting telescope 8300 provided with the second detector 8500 including the incremental optical encoder provided with the index; that is, the zero detecting index 8512 of the rotor 8510 must be detected by the index detecting unit of the optical detector 8530 by turning the sighting telescope 8300 one full turn. The sighting telescope 8300 needs to be turned by one full turn because the index 8512, in general, is invisible from outside. If the sighting telescope 8300 is turned quickly, the index detecting unit of the optical detector 8530 is unable to detect the zero detecting index 8512 of the rotor 8510 and hence measurement cannot be made. Accordingly, the sighting telescope 8300 must be turned at an appropriate turning speed below a certain turning speed, with the appropriate turning speed unavoidably being dependent on user intuition, which is very burdensome."}
-{"text": "1. Field of the Invention\nThe present invention relates to a vehicle seat in which a child seat is built into an occupant seat formed by a seat cushion and a seat back which is reclinable with respect to the seat cushion.\n2. Description of the Related Art\nThere are various types of conventional child seats. For example, there is a type in which the child seat is fixed to an ordinary occupant seat as a single unit. There is another type in which a child seat mounting/fixing member is built in an occupant seat so that the child seat is fixed to the occupant seat as a single unit. In yet another type, which will be referred to hereinafter as a \"built-in type\", a child seat is built in an occupant seat so that the vehicle seat can be changed between a state of use as an occupant seat and a state of use as a child seat as occasion demands.\nWith built-in type vehicle seats, there is no need to remove the child seat. Built-in type vehicle seats are therefore superior in that the only work involved therein is changing the seat from one state of use to the other (as an example, see Japanese Utility Model Application Laid-Open No. 2-23229).\nThe structure disclosed in the above-mentioned reference is as follows: in an ordinary occupant seat, a seat back and a connecting rod are swingable a predetermined amount with respect to a seat cushion. Vertical grooves are formed in a reverse surface of the seat back. A child restraining belt and a waist portion pad are accommodated in the vertical grooves. Accordingly, when the occupant seat is to be used, the vehicle seat is placed in its normal state. When the child seat is to be used, the seat back of the occupant seat is swung so that by using the seat cushion of the occupant seat as the seat cushion of the child seat as well, a child can sit facing the rear.\nHowever, in such a structure, the seat cushion of the occupant seat is used also as the child seat. Therefore, a drawback arises in that the angles of the seat surface are not structured so as to allow an optimal seat surface for a child. Further, in order to protect the sides of the child's body, it is necessary to form side wall portions. However, such side wall portions are unnecessary in an occupant seat. Namely, separate, independent cushion capabilities for an occupant seat and a child seat cannot be ensured."}
-{"text": "1. Field of the Invention\nThe invention relates to semiconductor wafer processing equipment and, more particularly, the invention relates to a method and apparatus for selectively marking a semiconductor wafer with identifiable markings.\n2. Description of the Background Art\nIdentifying defects on the surface of a wafer, such as particulate contaminants and other surface irregularities, is extremely important to integrated circuit manufacturing processes. To eliminate defect sources, defects are identified and analyzed to determine the source of the defect. Thereafter, corrective action can be taken to reduce or eliminate the defect source.\nTypically, the defect identification process is accomplished in two steps. First, a laser scanner device scans a wafer with a laser and analyses the backscatter of the laser to locate defects on the surface of the wafer. One such laser scanner device is a Tencor SurfScan 6200, manufactured by Tencor Instruments. Second, each defect located by the scanner is analyzed to identify the root cause of the defect. The most commonly employed analysis tool is a high magnification imaging system such as a scanning electron microscope (SEM). An SEM is used to identify the defect and/or the source of the defect by inspecting the defect at high magnification. Additionally, the SEM may be accompanied by instrumentation for performing chemical analysis of the defect. Such instrumentation includes an energy dispersive x-ray (EDX) detector. Other tools may include instrumentation for performing an Auger analysis, an atomic force microscope (AFM), a tunneling electron microscope (TEM), an optical spectrometer and the like.\nSince the SEM as well as other analysis tools use a high magnification (on the order of 200 to 2000 times), rapidly positioning the SEM at the defect location can be time consuming. Although the laser device provides defect coordinates that are accurate to approximately 300 microns, a substantial amount of time can be spent manually searching, for example, a 300 by 300 micron region with the SEM for a defect having dimensions of a few tenths of micron.\nAdditionally, when using a SEM (or other optical analysis tool) to analyze a \"bare\" wafer, i.e., a wafer having no surface features, the SEM has difficulty focusing on the wafer surface. Without an accurate focus, finding a small defect on the surface is nearly impossible.\nFurthermore, some defects are identified by laser scanning, but cannot be seen with a SEM. As such, the SEM operator may search for a defect for a long period of time until realizing the defect cannot be seen with the SEM.\nSome defect analysis tools, e.g., an atomic force microscope (AFM), are not readily useful when a large search area is used. Generally, the field of view for an AFM is approximately 10 .mu.m.sup.2, and the search area is as large as 300 .mu.m.sup.2. Unfortunately, an AFM requires approximately five minutes to obtain a 10 .mu.m.sup.2 image. As such, 900 images are required to cover a 300 .mu.m.sup.2 search area, requiring 4500 minutes to complete the search.\nTherefore, a need exists in the art for a method and apparatus that selectively marks a wafer proximate a defect such that the defect can be rapidly identified and examined with a SEM or other defect analysis system."}
-{"text": "Currently, front-facing cameras are configured on many terminal devices such as a mobile phone, a notebook computer, and a tablet computer, and are configured to shoot a photograph or a video of a target in front of the terminal devices.\nFor example, FIG. 1 shows a camera layout of an existing common terminal device such as a mobile phone. A rear-facing camera 101 is generally placed on the top or back of a screen 102 of the terminal device, and a front-facing camera 103 is generally placed on the top of the screen 102 of the mobile phone, so that the front-facing camera 103 can shoot a photograph or a video of a target in front of the terminal device.\nHowever, in an existing camera layout, relatively large space needs to be reserved at the top of a terminal device to place a front-facing camera, which causes a limited ratio of a screen to the entire terminal device."}
-{"text": "The present invention relates to energetic atomic beam generation, and particularly to electrically neutral molecular or atomic beam devices and methods for mass spectrometry and surface analysis, often called fast-atom bombardment mass spectrometry (FABMS).\nFABMS has several advantages over secondary-ion mass spectrometry (SIMS). The primary advantage is that FABMS allows the use of a liquid matrix, which simplifies sample preparation and maintains a reservoir of undamaged molecular sample species when subjected to an atomic beam under dynamic (intense particle flux) conditions. Secondly, the use of a neutral atom beam in the FABMS avoids the problem of floating the ion gun system above the accelerating voltage of the spectrometer. Finally, sample charge buildup is reduced with this neutral particle bombardment.\nMolecular SIMS is invaluable for the analysis and characterization of bulk solid surfaces, their films and molecular overlayers. The static (low particle flux) used in molecular SIMS is desirable to avoid damage to the molecular solid sample. Although conventional ion sources are capable of operation under such static conditions, attempts to charge neutralize the sample with an electron flood gun to avoid sample charge buildup are often ineffective, and they can result in bombardment-induced radiation damage and electron-stimulated desorption. The dynamic (intense particle) flux of conventional atomic beam sources quickly destroys the molecular overlayers or thin molecular film samples associated with the molecular SIMS method."}
-{"text": "1. Field of the Invention\nThe present invention relates to an antenna cover that protects an antenna of a plasma generating device.\n2. Description of the Related Art\nIn a plasma generating device, for example, plasma electrons are accelerated by an RF electric field generated in the vicinity of an antenna in a plasma chamber and plasma is generated. At this time, if the antenna is exposed to the plasma, the antenna is sputtered, which results in becoming the pollution cause of the plasma. For this reason, a member for protecting the antenna from the generated plasma is provided. For example, a partition plate is provided between the plasma chamber and an antenna chamber to protect the antenna.\nBecause the protection member of the antenna is exposed to the plasma, the protection member is damaged by the plasma. If the protection member is abraded under an influence of the plasma and a hole is formed in the protection member, the antenna is exposed to the plasma and the plasma is polluted. If the protection member is formed thick or the antenna or the protection member is arranged to be distant from a generation position of the plasma, abrasion of the protection member can be decreased. However, generation efficiency of the plasma is deteriorated."}
-{"text": "The present invention relates to an optical disk and a method of producing the same and more particularly to a stamper for molding an optical disk base highly compatible with commercially available CD (Compact Disk) players, a method of producing the stamper, a method of producing an optical disk base, a method of producing an optical disk, and an optical disk base, and an optical disk.\nIn parallel with the spread of optical disks, there is an increasing demand for the timely delivery of high quality optical disks to the market. Particularly, to enhance quantity production of optical disks, it is necessary to reduce a disk base molding cycle.\nTo produce an optical disk, a stamper formed with a transfer surface is positioned in one of a pair of mold parts forming a cavity therebetween. Molten resin is injected into the cavity and then cooled off. Subsequently, the mold parts are separated in order to remove the cooled resin. As a result, the transfer surface of the stamper is transferred to the resin, forming a recording surface.\nIt is a common practice with an optical disk to hold the mold parts at a temperature of about 200xc2x0 C. lower than the temperature of resin to be injected into the cavity. This promotes the cooling and solidification of the resin injected into the cavity. Such a mold temperature is determined by the tradeoff between transferability and an increase in the tact of a disk base molding cycle. Specifically, the mold temperature should be as low as possible for increasing the tact, but would degrade transferability if excessively low. On the other hand, a high mold temperature would enhance transferability, but would increase a period of time necessary for the resin to be cooled to a parting temperature and would thereby lower the yield of optical disks.\nJapanese Patent Laid-Open Publication Nos. 7-178774, 10-149587 and 6-259815 each propose to provide a mold or a stamper with a heat insulating ability so as to enhance both the transferability and the tact of the disk base forming cycle. Specifically, Laid-Open Publication No. 7-178774 teaches a heat insulating body removably positioned in a mold in such a manner as to face the rear of a stamper. Laid-Open Publication No. 10-149587 teaches a heat insulating ceramic layer formed on a mold in such a manner as to face the rear of a stamper. Further, Laid-Open Publication No. 6-259815 teaches a stamper whose front (transfer surface) is plated with Ni (nickel) containing 20% to 30% of polytetrafluoroethylene by electroless plating. Polytetrafluoroethylene has a grain size of 1.0 xcexcm or less. The resulting Ni film is 50 nm to 70 nm thick.\nHowever, none of the above conventional technologies can enhance both the transferability and the tact of a disk base molding cycle at a high level. Laid-Open Publication No. 6-259815 has a problem that the Ni film formed on the transfer surface of a stamper obstructs the fine patterning of the transfer surface. Laid-Open Publication No. 10-149587 has a problem that the mold itself must be redesigned or replaced, wasting existing molding equipment.\nSpin coating has customarily been used to coat a molded disk base with an organic pigment which forms a recording layer because spin coating is desirable from the easy process and low cost standpoint. While the thickness distribution of the recording layer can be controlled on the basis of coating conditions, it is difficult to control the distribution of the pigment in guide grooves. Specifically, to form the recording layer, a disk base is caused to spin such that a pigment solution sequentially spreads outward over the entire disk base due to a centrifugal force. However, the centrifugal force differs from one position to another position in the radial direction of the disk base. This, coupled with the fact that the solvent evaporates while spreading outward, causes the pigment to fill outer guide grooves more easily than inner guide grooves.\nIt follows that if the guide grooves of the disk base have a uniform configuration from the inner circumference to the outer circumference, the configuration of the guide grooves filled with the pigment differs from one position to another position in the radial direction. This scatters reflectance and tracking error and other signal characteristics and makes it difficult to produce constant quality, reliable optical disks. In addition, the resulting optical disks are not satisfactorily compatible with commercially available CD players.\nJapanese Patent Laid-Open Publication Nos. 5-198011 and 5-198012, for example, disclose implementations for correcting the above difference in configuration between the inner guide grooves and the outer guide grooves filled with the pigment. The implementations are such that the configuration (depth) of the guide grooves to be formed in a disk base or a stamper is intentionally varied beforehand. None of such implementations, however, gives consideration to the decrease in the fluidity of molten resin ascribable to temperature fall. Therefore, the implementations cannot realize desirable transferability alone when a high cycle is desired, aggravating the scattering of optical disks in signal characteristics.\nIt is therefore an object of the present invention to enhance both the transferability and the tact of a disk base molding cycle at the same time.\nIt is another object of the present invention to allow a transfer surface to be finely patterned.\nIt is yet another object of the present invention to make it needless for existing molding equipment to be redesigned or replaced.\nIt is a further object of the present invention to provide an optical disk sufficiently compatible with commercially available CD players by allowing guide grooves filled with a pigment by spin coating to have a substantially uniform configuration at any position in the radial direction.\nIn accordance with the present invention, a stamper for molding an optical disk base includes a transfer surface for molding the optical disk base, and a heat insulating material extending in parallel to, but not contacting, the transfer surface.\nAlso, in accordance with the present invention, a method of producing a stamper for molding an optical disk base includes the steps of electroforming on a photoresist master having a transfer surface pattern an Ni layer having a transfer surface to which the transfer surface pattern is transferred, forming a heat insulating layer on the Ni layer, and separating the photoresist master from the Ni layer.\nFurther, in accordance with the present invention, a method of producing an optical disk base includes the steps of injecting molten resin into a cavity formed by a pair of mold parts and accommodating a stamper having a transfer surface for molding the optical disk base and a heat insulating layer extending in parallel to, but not contacting, the transfer surface, and separating the pair of mold parts to thereby remove the resin cooled off.\nFurthermore, in accordance with the present invention, a method of producing an optical disk includes the steps of injecting molten resin into a cavity formed by a pair of mold parts and accommodating a stamper having a transfer surface for molding the optical disk base and a heat insulating layer extending in parallel to, but not contacting, the transfer surface, separating the pair of mold parts to thereby remove the resin cooled off, coating a transfer surface of the resin with a recording material to thereby form a light absorption layer, forming a reflection film on the light absorption film, and forming a protection film on the reflection film.\nMoreover, in a method of producing an optical disk base, a heat insulating material is positioned beneath a recording area formed on the surface of a stamper for molding an optical disk."}
-{"text": "Present apparatus and methods for changing out carbon in an industrial carbon filter system, such as in a chemical process plant, include 50 to 2000 pound carbon canisters having openings at the top for receiving fresh carbon. Releasable attachment parts affixed to the canisters, such as nipples and flanges, are provided for connecting the canister into a process stream line for filtering the stream. A stream inlet nipple is usually provided near the top of the canister and a stream outlet nipple near the bottom.\nThe plant typically provides a roll-off box or the like for collecting spent carbon and sending it off to be revitalized and then re-used. The plant also usually provides a fork lift for moving canisters around the site. Typically a service provider is retained by the plant, however, to change out spent carbon and to refill the canisters with fresh carbon, upon demand and need. Spent carbon is shipped off-site to be revitalized. This spent carbon may be wet or dry.\nCurrent service provider systems and methods typically include a truck, manned by at least two people, the truck having a hopper and vacuum apparatus. At least two people are required to set up the hopper and run the vacuum operation. In operation, in the prior art, canisters with spent carbon filter material are disconnected from their process streams and brought to a central site, usually proximate the location of a roll-off box or the like at a plant facility. A service provider vacuum truck is driven to the site. The at least two service people set up the hopper as the well as the vacuuming equipment. The upper lid of the canisters are opened, each in turn, and spent carbon is vacuumed out from the canister and into the provider's hopper. When the hopper is filled, (a hopper typically holds about 1500 pounds of spent carbon,) vacuuming is stopped and the hopper is frequently moved by the fork lift to a roll-off box, dumped into the roll-off box and returned to the proximity of the truck for further vacuuming. In a less favored system, because it entails higher risks, the hopper dumps its contents into one or more spent carbon container bags at the canister site. The bags are then tied off and moved to a location for collecting and shipping off site. (The possibility of the deterioration of the bags and resulting contamination of the site prior to shipping causes this option to be slightly less favored.)\nThe instant invention comprises an improvement to these known prior art systems and methods. The instant invention, at least in certain key applications, has been tested to determine that it is more cost effective and requires less equipment and fewer people and takes less time.\nIn preferred key embodiments of the instant invention, a canister is provided with two valved orifices. While one orifice may always be provided on top of the canister for re-filling the canister with fresh carbon, preferably a second valved orifice is provided by the instant invention on lower parts of the canister or on the bottom. This second orifice is preferably situated on the bottom and preferably includes a ball valve or the like for controlling an approximately four to six inch diameter opening. The same roll-off box and forklift provided by the plant, as discussed above, can be used in the instant invention. As will be shown, however, the instant invention can also be utilized with existing canisters.\nIn key preferred embodiments of the instant invention, a truck arrives with one person to the site of the canisters. (The truck may include vibration equipment for difficult canisters.) In this preferred embodiment one person lifts, with the fork lift, each canister up and over the top of the plant's roll-off box or the like. (Alternate embodiments are also disclosed below.) The service person opens a lower orifice on the bottom of the canister using a ball valve, and the spent carbon drops by gravity into the roll-off box. Alternately, the service person might pivot the canister, to invert it over the roll-off box, such that the spent carbon dumps from the top orifice.\nThe instant inventor's tests indicate that while it takes two people 8-9 minutes to vacuum a thousand pounds of spent carbon from a canister, it only takes one person two minutes or less to dump a thousand pounds of spent carbon into a roll-off box. (Note: the canister could also be dumped into saturated carbon containers or bags instead of a roll-off box. Again, use of such bags is slightly disfavored, as mentioned above, for incurring higher potential contamination hazards.)\nThe instant invention eliminates the need for vacuuming equipment, for the second or more service person and for the intermediate step of filling and emptying the service provider's hopper.\nIn the odd situation when spent carbon might resist dropping out of a canister by gravity, tests have indicated that the techniques of shaking, scraping and/or vibration can be successfully employed. The inventor's tests have demonstrated that a forklift can affectively shake a canister to cause sections of spent wet carbon that failed to fall by gravity, to break off and fall out of a canister. The inside walls of a canister can also be scraped, or squeegeed, by the service person. In addition, a vibrator, such as a vibrator used with a cement mixer, can be connected or inserted in the canister, possibly with a hand held prong inserted inside of the canister, to help ensure that all spent carbon is dislodged.\nAgain, preferred novel canisters of key embodiments of the instant invention have a second bottom valved orifice, which opens and closes. Although these canisters should be more expensive than the canisters of the prior art, that cost should be more than recovered by the enhanced speed and lessened cost of dumping rather than vacuuming, and also by the reduced cost of manpower and equipment."}
-{"text": "During seismic surveying, a seismic source is repeatedly actuated and seismic receivers, such as hydrophones in marine seismic surveying, receive energy direct from the sources and reflected from various boundaries or interfaces. In the case of marine seismic surveying, energy propagates into the earth and is reflected back to the hydrophones from subterranean boundaries or interfaces, for instance between strata of different types. The recorded seismic data can be processed to reveal information about the structure of the earth in the area being surveyed. However, such reflections are contaminated by other reflection paths. For instance, energy from the sources is reflected at the sea surface directly to the hydrophones. Also, energy can be reflected more than once between the sources and the receivers. Such multiple reflections can take place within the earth. Also, energy initially traveling downwards from the sources can be reflected upwardly and then downwardly again from the surface of the sea before arriving at the hydrophones. Reflections of this type are referred to as \"free-surface multiple reflections\". Free-surface multiple reflections can be classified according to their order, which is equal to the number of reflections from the free-surface. Thus, first order free-surface reflections comprise energy initially traveling downwardly from the sources (as opposed to \"ghosting\" where energy travels upwardly and is reflected from the free-surface), is reflected upwardly from the sea bed or a boundary below the sea bed, and is then reflected downwardly from the free-surface to the hydrophones. Second order free-surface multiple reflections undergo two downward reflections from the sea-surface before being detected by the hydrophones, and so on.\nThe invention makes use of a general concept of multiple attenuation using the inverse scattering or scattering Born series. For this general concept reference is made to the published International patent application WO 95/10787 and the British patent application No. 9426255.7 (published as GB-A-2296567).\nFrom those documents it is known how to derive the scattering Born series for removing free surface multiples: EQU D.sub.P =D.sub.0 +AD.sub.1 +A.sup.2 D.sub.2 +A.sup.3 D.sub.3 + . . . ,[1]\nwherein D.sub.P denotes the signals without free-surface multiples and A denotes the inverse of the source signature. The first term of the scattering series, D.sub.0, is the actual seismic data, the second term, D.sub.1, removes first order free-surface multiples, the following term, D.sub.2, removes second order free-surface multiples, and so forth.\nIt is an object of the present invention to improve and utilize the above method for sea bottom seismic acquisition operations."}
-{"text": "The present invention relates to disposable containers for hospital sharps and infectious waste, and pertains particularly to a secure disposable container and alternate closure assemblies for stand alone and for insert containers.\nHospitals and medical clinics use great quantities of sharps, such as needles, syringes, surgical blades, and the like, that are disposed of rather than cleaned and reused. It is necessary that the sharps be disposed of in a manner that prevents them being reused without sterilization. In particular, it is necessary to keep them from falling into the hands of those, such as intravenous drug users and the like, who are likely to use them without proper sterilization.\nNumerous containers have been developed in recent years, which are reasonably secure and disposable for receiving and disposing of hospital sharps, wastes and the like. Many of these containers are disposable inserts that mount into a security housing which is lockable and is securely mounted on a wall or other support structure. Many of these insert containers however have an inadequate security closure built into it, and are out dated so that they do not provide adequate security against pilfering of used syringes and the like therefrom. While improved containers have been developed which cannot readily be reopened and articles cannot be easily removed therefrom, such containers must be kept in a secure place or securely mounted to non-removable structure to prevent unauthorized removal.\nMany hospitals and clinics that have the older secure housings with insert type disposable containers do not wish to invest in new housings. It is therefore necessary to develop improved insert type disposable containers for these housings with improved closures for the combination. It is also desirable that new stand alone units be available.\nAn example of a prior art disposable insert type container and mounting housing is disclosed in U.S. Pat. No. 4,715,498 of Hanifl dated Dec. 29, 1987. This patent discloses a sharps container having a closure door that provides limited access to the interior of the container for mounting in a protective housing. However, an improved closure is desirable for such containers.\nExamples of stand alone type containers are disclosed in U.S. Pat. No. 4,736,860, granted Apr. 12, 1988 to Bemis, and entitled SHARPS DISPOSAL CONTAINER, and in U.S. Pat. No. 4,828,107, granted May 9, 1989 to Spencer and entitled DISPOSABLE CONTAINER FOR SYRINGES. It is desirable that simpler more effective disposable containers be available.\nIn the manufacture of multiple related or similar products, it is desirable that as many parts be interchangeable as possible in order to reduce manufacturing costs. It is also desirable that different variations of the product be available to meet different conditions and needs. It is, therefore, desirable that an improved securable disposable container combination be available to provide both insert and stand alone containers."}
-{"text": "A single crystal used as a substrate of semiconductor devices, for example, a silicon single crystal, is mainly produced by Czochralski method (hereafter abbreviated to CZ method).\nWhen producing a single crystal by CZ method, for example, a single crystal production apparatus 10 as shown FIG. 2 is used to produce the single crystal. This single crystal production apparatus 10 has a member for containing and melting a raw material polycrystal such as silicon, heat insulation members to insulate heat, and etc. They are installed in a main chamber 11. A pulling chamber 12 extending upwardly is continuously provided from a ceiling portion of the main chamber 11, and a mechanism for pulling a single crystal 13 by a wire 14 (not shown) is provided above it.\nIn the main chamber 11, a quartz crucible 16 for containing a melted raw material melt 15 and a graphite crucible 17 supporting the quartz crucible 16 are provided, and these crucibles 16 and 17 are supported by a shaft 18 so that they can be rotated and moved upwardly or downwardly by a driving mechanism (not shown). To compensate for decline of the melt level of the raw material melt 15 caused by pulling of the single crystal 13, the driving mechanism for the crucibles 16 and 17 is designed to rise the crucibles 16 and 17 as much as the melt level declines.\nAnd, a graphite heater 19 for melting the raw material is provided so as to surround the crucibles 16 and 17. A heat insulating member 20 is provided outside the graphite heater 19 so as to surround it in order to prevent that the heat from the graphite heater 19 is directly radiated on the main chamber 11.\nMoreover, a graphite cylinder 23 is provided above the crucibles, and a heat insulating material 24 is provided on the outside of the lower end of the graphite cylinder 23 so as to oppose to the raw material melt 15 so that the heat radiation from the melt surface is intercepted and the temperature of the raw material melt surface is kept.\nA raw material lump is put in the quartz crucible 16 provided in the single crystal production apparatus as described above, the crucible 16 is heated by the graphite heater 19 as described above to melt the raw material lump in the quartz crucible 16. A seed crystal 22 fixed by a seed holder 21 connected with the lower end of the wire 14 is immersed into the raw material melt 15 melted from the raw material lump as described above. Thereafter, the single crystal 13 having a desired diameter and quality is grown under the seed crystal 22 by rotating and pulling the seed crystal 22. In this case, after bringing the seed crystal 22 into contact with the raw material melt 15, so-called necking, once forming the neck portion by narrowing the diameter to about 3 mm, is performed, and then, a dislocation-free crystal is pulled by spreading to a desired diameter.\nThe silicon single crystal produced by the CZ method as described above is mainly used for producing a semiconductor devices. In these years, with higher integration of semiconductor devices, an element becomes finer. The problem of Grown-in crystal defects introduced during the crystal growth has become more significant as an element has become finer.\nHereafter, the Grown-in crystal defects will be described (see, FIG. 4).\nIn the silicon single crystal, if the crystal growth rate is relatively high, Grown-in defects such as FPD (Flow Pattern Defect) considered to be originated from voids aggregated of void-type point defects, are present at high density over the region in a crystal radial direction. The region where the defects are present is called V (Vacancy) region. Moreover, if the growth rate is lower, OSF (Oxidation Induced Stacking Fault) region is generated as in a ring shape from the periphery of the crystal with lowering of the growth rate. And defects such as LSEPD (Large Secco Etch Pit Defect) and LFPD (Large Flow Pattern Defect) considered to be originated from dislocation loops aggregated of interstitial silicons in the outside of the ring are present at low density, and the region where the defects are present is called I (Interstitial) region. With further lowering of the growth rate, OSF ring is contracted to disappear in the center of a wafer, and the whole plane becomes I region.\nIn these years, in the outside of the OSF ring between V region and I region, a region has been discovered where neither such as FPD originated from voids nor such as LSEPD and LFPD originated from interstitial silicons are present. The region is called N (Neutral) region. Moreover, the N region is further categorized as follows. There are Nv region (a region where more voids are present) next to the outside of the OSF ring and Ni region (a region where more interstitial silicons are present) next to I region. In the Nv region, amount of precipitated oxygen is rich when thermal oxidation treatment is performed, and in the Ni region, amount of precipitated oxygen is little.\nThe Grown-in defects are considered to be determined the introduced amount by the parameter of F/G which is ratio of the pulling rate (F) and an average value of a temperature gradient in the crystal along a pulling axis from the melting point of silicon to 1400\u00b0 C. (G) (See, for example, V. V. Voronkov, Journal of Crystal Growth, 59 (1982) 625-643). Namely, if the pulling rate and the temperature gradient are controlled so that F/G is constant, a single crystal can be pulled in a desired defect region or in a desired defect-free region (See, for example, Japanese Patent Laid-open (Kokai) Publication No. 2000-178099).\nTherefore, it is conventionally necessary to pull a single crystal by controlling the pulling rate etc. for N region so as to obtain a single crystal of defect-free region. It is difficult to control the rate because the single crystal of N region can be grown in a range of a relatively limited pulling rate. Thus, productivity and yield of the single crystal become low. Accordingly, a method of expanding a pulling rate range where a defect-free region can be obtained has been required so as to produce it more easily."}
-{"text": "1. Field of the Invention\nThe present invention relates to protective knee and elbow protectors employed by athletes and the like, particularly adapted to prevent injuries to the knee and the proximal anatomical structure associated therewith.\n2. Description of the Prior Art\nVarious knee and elbow protector devices are known in the art. It has been found that by placing a protector element atop the kneecap, injuries sustained in physical activity are reduced. Further, it has been determined by using such devices, the incidence of aggravating a current injury is also reduced. As such, knee and elbow protector devices are beneficial in the sport and athletic arts.\nExamples of knee and elbow protector protectors and other protector devices for use during athletic activity are described in the following U.S. Patents:\nU.S. Pat. No. 2,641,761 to Schultz discloses a knee brace or stabilizer. PA1 U.S. Pat. No. 3,044,075 to Rawlings discloses a protective device to protect the wearer at impact points. PA1 U.S. Pat. No. 3,322,873 to Hitchcock discloses a resilient body protector for use over various portions of the body. PA1 U.S. Pat. No. 3,670,725 to Gaylord discloses a protective pad to be employed generally with jointed portions of the body. This pad would be placed atop the jointed element to prevent debilitating bedsores and ulceration. PA1 U.S. Pat. No. 4,250,570 to Barlow discloses a protective knee support. PA1 U.S. Pat. No. 4,287,885 to to Applegate discloses a knee brace with a resilient pad for surrounding the patella. There exists no central pad, merely a central opening to receive the extended portion of the patella therethrough. PA1 U.S. Pat. No. 4,354,280 Hayes discloses a contact sport joint protector. PA1 U.S. Pat. No. 4,484,361 to Leighton discloses a knee and elbow pad and a method of manufacturing the same. PA1 U.S. Pat. No. 4,573,216 to Wortbeg discloses an impact dissipater to protect bones from impacts or falls. This device is primarily designed to be used proximal the bony region of the hip portion of the body. PA1 U.S. Pat. No. 4,700,698 to Aleylein discloses a knee orthosis. This device includes a right and a left lateral extension which is integral to the elastic cylindrical element. The device does not include a central foam element nor does it include an upper or lower lateral extension. PA1 U.S. Pat. No. 4,723,322 to Shelby discloses a knee pad. PA1 U.S. Pat. No. 4,796,303 Atwater discloses a knee protector pad. PA1 U.S. Pat. No. 4,870,956 to Fatool discloses a knee brace. PA1 U.S. Pat. No. 5,077,837 to Meistrell discloses a knee or elbow protector. PA1 U.S. Pat. No. 5,105,473 to Valtakwi discloses an athletic outfit having pockets for protective pads. PA1 U.S. Pat. No. 5,306,229 to Brandt discloses an elastic articular bandage. PA1 U.S. Pat. No. 5,411,037 to Hess discloses a knee-joint bandage. PA1 U.S. Pat. No. 5,474,524 to Carg discloses a joint support. PA1 U.S. Pat. No. 5,680,657 to Valtakwi discloses protective pads for insertion into pockets to protect from impact while playing sports. PA1 U.S. Pat. No. 5,685,021 to Tsujino discloses protective head gear for wrestlers containing a gel cushioning material. PA1 U.S. Pat. No. 5,695,452 to Grim discloses a formed resilient orthopedic pad. PA1 U.S. Pat. No. 5,769,809 to Witzell discloses a below the joint amputation limb protector. PA1 U.S. Pat. No. 327,961 to Decker, U.S. Pat. No. 341,005 to Pratt, U.S. Pat. No. 341,229 to Tsuiino, U.S. Pat. Nos. 360,284, 360,285, 360,286 and 360,287 to Pafet, U.S. Pat. No. 363,147 to Chen, U.S. Pat. Nos. 394,332 and 394,333 to Tsujino.\nThe following US design patents include pads and protectors for knees and other joints:\nAll of the forementioned patents are directed to protecting the knee or other body part from a direct impact. The instant invention is directed to a knee and elbow protector having a specific structure that protects the knee from an oblique, non-direct, or proximal impact. None of these references, either singly or in combination, disclose the instant invention and the benefits derived therefrom."}
-{"text": "1. Technical Field\nThe invention relates to a multiple-ratio transmission mechanism in a powertrain for an automotive vehicle and to a control strategy for achieving smooth engagement and release of friction torque establishing elements during a transmission upshift event.\n2. Background\nIn a geared automatic transmission in an automotive vehicle powertrain having an engine or other torque source, a ratio change may be made from a so-called low ratio to a so-called higher ratio when a friction torque establishing element, such as a clutch or brake, is engaged in synchronism with disengagement of a companion friction torque establishing element. This is referred to as a ratio upshift. The friction torque establishing elements involved in the upshift may be referred to as an oncoming clutch or brake and an off-going clutch or brake. The upshift event is characterized by a preparatory phase, a torque phase and an inertia phase as the vehicle accelerates from a standing start.\nIn a conventional automatic transmission in a vehicle powertrain, the oncoming clutch torque capacity is controlled to increase from a low value during the torque phase. Simultaneous engagement of one clutch or brake and release of another results in a momentary activation of two torque flow paths through the gearing, causing a gear tie-up in which transmission output shaft torque decreases momentarily. This condition may be referred to as a \u201ctorque hole\u201d. It occurs before the off-going clutch totally disengages.\nFriction elements, such as disc clutches, band brakes and disc brakes, typically are actuated hydraulically under the control of a transmission control module, which disengages an off-going friction clutch or brake while simultaneously engaging an oncoming friction clutch or brake during an upshift in order to lower speed ratio. For purposes of the present description of the invention, the clutch and the brake will be referred to as friction elements.\nDuring the preparatory phase, an automatic transmission control reduces off-going friction element torque capacity to prepare it for release as an actuator for the oncoming friction element is adjusted to prepare for its engagement. During the torque phase, the controller increases oncoming friction element torque capacity, which causes torque transmitted through the off-going friction element to drop quickly due to the transient gear tie up.\nAs torque is transmitted through the off-going friction element deceases, the automatic transmission output shaft torque drops, which causes the so-called torque hole. This is perceived by a vehicle occupant as an unpleasant shift shock. The inertia phase begins when the off-going clutch is released with no significant torque capacity."}
-{"text": "As fixed type constant velocity universal joints, there have been publicly known joints of a six-ball Rzeppa type (BJ), a six-ball undercut-free type (UJ), an eight-ball Rzeppa type (EBJ), an eight-ball undercut-free type (EUJ), and a track groove crossing type in which paired track grooves cross each other (see, for example, Patent Literature 1). Those joints are used as appropriate in accordance with purposes, required characteristics, and the like.\nReferring to FIG. 14a and FIG. 14b, description is given of an example of the fixed type constant velocity universal joint of the track groove crossing type (hereinafter referred to as \u201ccrossing fixed type CVJ\u201d). FIG. 14a is a vertical sectional view (sectional view taken along the line A-O-B in FIG. 14b) of a state in which the crossing fixed type CVJ for a propeller shaft forms an operating angle of 0\u00b0, and FIG. 14b is a front view of the crossing fixed type CVJ. The constant velocity universal joint 100 includes an outer joint member 102, an inner joint member 103, balls 104, and a cage 105. Eight arc-shaped track grooves 107 are formed in a spherical inner peripheral surface 106 of the outer joint member 102. The track grooves 107 are formed so that planes including ball raceway center lines x of the track grooves 107 are inclined with respect to a joint axial line n-n and the track grooves 107 are adjacent to each other in a circumferential direction with their inclination directions opposite to each other. Eight arc-shaped track grooves 109 are formed in a spherical outer peripheral surface 108 of the inner joint member 103. The track grooves 109 are formed so as to be mirror-image symmetrical with the paired track grooves 107 of the outer joint member 102 with respect to a plane P including a joint center O at the operating angle of 0\u00b0. That is, the inner joint member 103 is assembled to an inner periphery of the outer joint member 102 so that the paired track grooves 107 and 109 cross each other.\nAs illustrated in FIG. 14a, curvature centers of the track grooves 107 of the outer joint member 102 and the track grooves 109 of the inner joint member 103 are each positioned at the joint center O. The balls 104 are interposed in crossing portions between the paired track grooves 107 and 109, and are received and held in pocket portions 105a of the cage 105 arranged between the spherical inner peripheral surface 106 of the outer joint member 102 and the spherical outer peripheral surface 108 of the inner joint member 103. Curvature centers of a spherical outer peripheral surface 111 and a spherical inner peripheral surface 112 of the cage 105 are each positioned at the joint center O. In the constant velocity universal joint 100, the balls 104 are interposed in the crossing portions between the paired track grooves 107 and 109. Therefore, when the joint forms an operating angle, the balls 104 are always guided in a plane bisecting an angle formed between axial lines of the outer joint member 102 and the inner joint member 103. With this, rotational torque is transmitted at a constant velocity between the two axes.\nFurther, the track grooves 107 of the outer joint member 102 and the track grooves 109 of the inner joint member 103 are adjacent to each other in the circumferential direction with their inclination directions opposite to each other. Therefore, when both the joint members 102 and 103 rotate relative to each other at the operating angle of 0\u00b0, forces in the opposite directions are applied from the balls 104 to the pocket portions 105a of the cage 105 that are adjacent to each other in the circumferential direction. With this, the cage 105 is stabilized at the position of the joint center O. Thus, a contact force between the spherical outer peripheral surface 111 of the cage 105 and the spherical inner peripheral surface 106 of the outer joint member 102, and a contact force between the spherical inner peripheral surface 112 of the cage 105 and the spherical outer peripheral surface 108 of the inner joint member 103 are suppressed. As a result, it is possible to attain the constant velocity universal joint 100 that is capable of suppressing torque loss and heat generation, and is excellent in durability. Thus, with use of the constant velocity universal joint 100, it is possible to attain a propeller shaft that is suppressed in torque loss and heat generation, and is enhanced in efficiency.\nOperability of fixed type constant velocity universal joints including the constant velocity universal joint 100 described above, in which rotational torque is transmitted between both the joint members through intermediation of the balls, is secured by setting wedge angles of ball tracks formed by the paired track grooves (by holding the balls with the paired track grooves). In the constant velocity universal joint 100 of the track groove crossing type, as described above, when wedge angles, which are formed so that a force of pressing the cage 105 to one axial side (for example, right side in FIG. 14a) is applied from the balls 104 at the operating angle of 0\u00b0, are defined as positive wedge angles, and wedge angles, which are formed so that a force of pressing the cage 105 to the other axial side is applied from the balls 104 at the operating angle of 0\u00b0, are defined as negative wedge angles, ball tracks that form the positive wedge angles and ball tracks that form the negative wedge angles are alternately formed in the circumferential direction. With this, desired joint performance is secured. In the constant velocity universal joint 100 described above, the curvature centers of the track grooves 107 and 109 are each positioned at the joint center O, and hence sizes and opening directions of the wedge angles are determined in accordance with angles (inclination angles) that are formed by the track grooves 107 and 109 (by planes including the ball raceway center lines thereof) with respect to the joint axial line n-n, contact angles of the balls 104 with respect to the track grooves 107 and 109, and operating angles."}
-{"text": "1. Field of the Invention\nThe present invention is related to methods of protection against malware located on web resources and, in particular, to malware scans of web resources and identification of malware components on web resources.\n2. Description of the Related Art\nDetection of viruses and malware has been a concern throughout the era of the personal computer. With the growth of communication networks such as the Internet and increasing interchange of data, including the rapid growth in the use of e-mail for communications, the infection of computers and networks through communications or file exchanges is an increasingly significant consideration. Infections take various forms, but are typically related to computer viruses, Trojan programs, or other forms of malicious code (i.e., malware).\nRecent incidents of e-mail mediated virus attacks have been dramatic both for the speed of propagation and for the extent of damage, with Internet service providers (ISPs) and companies suffering service problems and a loss of e-mail and networking capability. In many instances, attempts to adequately prevent file exchange or e-mail mediated infections significantly inconvenience computer users.\nPopularization of web services increases malware-related threats to web clients. With the development of web-technologies, such as AJAX, JAVA, PHP, FLASH, etc., web sites become more accessible, with more functionality with various media content, which enhances their appeal and popularity among users. With the introduction of methods of designing web-systems, known as Web 2.0, users have been enabled to not only receive information from the Internet, but also fill out web sites with their own data.\nMany Internet users have their own blogs, pages within social networks, where they can share messages or data with other users. In addition to text messaging, modern technologies support transfer of media files, such as photos, video files, interactive documents, animations and applications.\nActive participation of users in creation and modification of a web site leads to a rapid development and changes of the site content. A typical example of such a site is a news ticker, where new articles with links to news reports, photos or videos appear at intervals of a minute and sometimes even more often. Another example is a web forum, where a number of users can be in excess of several thousand, and where new messages appear every second.\nIncreasing popularity of web resources makes this environment more attractive to hackers and virus writers, who spread malicious programs over the Internet. A malicious script, such as iframe, exploits, etc. can be added to the files uploaded to the site. Thus, file exchanges, forums, blogs, web-interfaces of mail servers and any other resources can be infected.\nThe infection can be perpetrated through any interface, such as a usual web site form to be filled out by a user (in this case, the infection is performed mostly manually by an insider or an intruder who stole passwords). A perpetrator can exploit vulnerability in the content management system or file access the site via FTP. The list of resources is not limited to the http transfer protocol but also includes ftp-resources and other servers.\nTypically, ftp-servers are used as tools for remote administration of the sites, including editing and uploading scripts. However, a connection to the ftp-server provides the ability to test all scenarios of site, and not just those that are executed on the user's computer. Access the site files via FTP allows a security system to analyze the source scripts and original pages, while access via http allows to examining only the result of the script and/or the result of processing of the web site files by a web server.\nAn unauthorized access to a user account and uploading on behalf of this user, a malicious program could lead to a rapid spread of malware because of the credibility of that person among other users. When dealing with mundane sites, users' attention may fade over time due to the usual circumstances. Because of this, incidents of exploitations of social networks and phishing are increasing.\nIn addition to public Internet services, there are some local sites (corporate, local area networks, user group sites, etc.) that are not accessible from outside of their network. This imposes significant problems in malware scanning of a given resource by anti-virus services or laboratories.\nA typical scheme of interaction between a user computer system and a Web resource is illustrated in FIG. 1. Web-server is an application that performs the functions of the server (i.e., a computer system) on which the application is implemented. In addition to web-servers 110 or ftp servers 120, other server applications can be installed. For example, mail servers, proxy servers, IRC-server, etc. In order to interface with these servers, client (i.e., a user) computers 140 must have special applications installed, such as, a browser 130, a file manager 150, etc.\nClient 140 access the web server 100 using a URL address of a desired web page or other resource. Each server application and a client interact, using different protocols. The primary protocols for data transfer between the client 140 and server 100 are HTTP (HTTPS), (S) FTP, POP3, SMTP, IMAP4, etc.\nMost servers support authentication in the following manner: an authorization data is transferred from the client 140 to the server 100, where the authorization takes place. Then, the data is transmitted from the server 100 to the client 140 based on the rights granted to the client 140.\nPersonalization makes it possible to make a web resource unique to each user. Authorization is used to distinguish users, their rights and available data. User authorization is carried out by the user via a client application, for example, by filling the forms in the browser. Very often these applications allow preserving the identity of the user. Thus, each user can customize his own interface, to restrict access to his personal page or email, to identify the displayed pages by sections or topics.\nIn order to provide security while using Internet or on the local network, downloadable content needs to be checked. A downloaded page may contain viruses, trojans, adware, spam, exploits (i.e., HTML code, links to image or \u201c.pdf\u201d file, which has special modifications that result in browser errors and execution of some potentially harmful code) for applications (such as for example, PDF-reader, web browser, media player, flash-player, etc.)\nA system of protection of personal computers can include a file anti-virus (AV), network screen, a firewall, a special protection against network attacks, Web anti-virus as well as remote security means of AV vendor companies. Currently AV technologies are heavily developed and includes many different methods and systems that implement both heuristic and signature analysis.\nThe signature type web analysis includes:\nassembling black list of pages (URL-blacklist);\nassembling white lists of trusted (i.e., clean) applications/components; and\nstoring a collection of malware components.\nThe heuristic analysis typically includes any of:\nemulation of executable programs/components;\nemulation of executable scripts;\nvirtualization of execution environment; and\ncontrol and analysis of application activity.\nA firewall or a network screen is necessary for monitoring and filtering network packets at different levels of the OSI model in accordance with assigned rules for network connections: filtering based on static rules and filtering with tracking executed applications and controlling logic and algorithms of the relevant protocols.\nA system for protection against network attacks is typically launched at system startup and monitors incoming traffic activity for patterns typical of network attacks. If an attempt to attack a computer is detected, the system blocks any network activity for the attacker to a protected computer.\nA conventional web anti-virus (web AV) intercepts and blocks execution of a malicious script on a web site if it poses a threat. Strict control is also imposed over all HTTP-traffic. Web AV also analyzes web resources for phishing scams and filters banners and pop-ups.\nAnother line of defense can be a local security server that analyzes the situation within a local network, scans the local traffic and analyzes the network activity of computers. There are also web services that allow the verification of the Internet resource or files. The user downloads a file or enters a URL of the resource and the entire malware test takes place on the web-service of an AV company.\nIn the case of local AV remedies, a web resource is checked when it is downloaded, in other words, when the client application goes to the corresponding URL-address. It is important to note that in this case, the user is authenticated on the web resource and downloads content, as defined for his user account. The content can be in form of scripts, links, articles, messages, reports, letters from trusted users, etc.\nFIG. 2 illustrates a conventional system for checking a web resource 200 for presence of malware or links to infected resources. A connection with the resource 200 is established by an application 220 of a client 210. The client application 210 in this case can be a regular browser, a file manager or another application that interacts with a server 200 via data transfer protocol 230.\nThe client application 220 transmits to the server 200 user identifiers 250. Depending on authentication scheme 260, the identifiers 250 can represent logins, passwords, session keys, cookie-files, special protocol headers, network or physical addresses of the computer, biometric data, certificates, etc. A request to the server 200 by the client-application 220 is processed by a server application 240.\nAn authorization 260 is performed, and depending on its results, the server 200 opens a document or generates a web page, an ftp-page or other data representation based on the data type provided by the server 200. Security of data transmitted between the client 210 and the server 200 is provided by a security module 270. This conventional system can be implemented on a personal computer as well as on a web server of an AV company. Typically, the relevance of the AV databases and the effectiveness of heuristic analysis on the server side can be higher than on a personal computer.\nHowever, the principal difference in this case is the data presented by a web resource 200 being tested for malware presence based on the user identification data 250. When the AV check takes place on the server, the result of the authorization is either denial of access or a grant of the rights of the guest account. Thus, the outcome of an AV check, even at equal technological capabilities on the server and the user computer, will be significantly different.\nYet another shortcoming of conventional AV systems is restrictions imposed upon AV web-based scanners by routing rules. If a web resource is part of the network and has no external network address, it is not available outside the network, and it can only be verified by using other security tools installed on a computer system connected to this network.\nMany malware creators and hackers are aware of online scanners and knowingly block access of the scanner applications to a web resource where they have planted malware components. This makes comprehensive malware scans difficult. Furthermore, in the conventional systems, a web page cannot be scanned until it is downloaded onto a user computer. A typical web resource can contain several thousand pages. AV checking all of the web pages using a conventional system (as illustrated in FIG. 2), requires a lot of time and resources.\nThe situation is further complicated when several different web resources need to be periodically scanned. Statistics indicate that the majority of malware components and links to infected pages are located on the main (i.e., home) pages of web resources or on the first pages of sub-sections. This also complicates AV checking process, since each server often uses different identification parameters. In order to check a list of web resources, an AV application needs to have access authorization that also complicates scanning web sites for malware.\nAccordingly, there is a need in the art for a comprehensive malware scanning system that can effectively check the web resources with a minimal overhead and costs."}
-{"text": "This invention relates to the art of beauty parlor equipment and, more particularly, to a drainage tray for a shampoo bowl to facilitate washing and rinsing of a patron's hair.\nIt is of course well known that beauty parlors have shampoo bowls over which a patron's head is disposed during shampooing and rinsing of the patron's hair. Most often, the bowl has an arcuate section in a front wall thereof which cradles the patron's neck, and the rear end of the bowl includes water supply valving and a spray head for rinsing shampoo from a patron's hair following washing thereof. The shampoo bowl further includes a drain in the bottom wall thereof towards the rear of the bowl through which rinse water and shampoo drains during a rinsing operation.\nIt is not unusual for a patron to have shoulder length or even waist length hair, and the shampooing and rinsing thereof is not only difficult but potentially unsanitary to some extent. In this respect, for example, the hair of a patron with shoulder to waist length hair lies on the bottom of the shampoo bowl which makes it difficult for rinse water to penetrate the hair without the beautician having to lift the hair off of the bottom of the shampoo bowl with one hand while trying to rinse the hair with the other hand. Moreover, if the patron's hair is left on the bottom of the shampoo bowl a residue often occurs on the ends of the patron's hair, and the hair is also subject to any contaminants which might be left in the bowl from a previous patron and/or from solutions used by the beautician in connection with washing and rinsing a patron's hair. Still further, if the patron's hair is of waist length or thereabout, it can easily flow into the drain opening which not only can cause clogging of the latter but also adds to the potential unsanitary condition of the hair. Again, to avoid the latter, a beautician must support the patron's hair in one hand and try and rinse it with the other which is time consuming and difficult as well as inefficient."}
-{"text": "1. Field of the Invention\nThe present invention relates to a polymer film exhibiting high retardation along thickness direction (Rth); a phase difference film, a polarizing plate, and a liquid crystal display device using the polymer film; and an Rth inducing agent.\n2. Description of the Related Art\nThere have been known techniques of adding various additives to a polymer film, aiming at imparting thereto various functions.\nAn exemplary proposal has been made on a polymer film improved in ultraviolet absorption, by means of addition of a merocyanine-based compound, as a ultraviolet absorber, to the polymer film (i.e., JP-A-H08-239509). JP-A-2009-67973 discloses a polymer material containing a merocyanine-based compound as an ultraviolet absorber, a mold product manufactured by using the polymer material, and also a coated ultraviolet absorber layer. Other proposals have been made on use of a merocyanine-based compound as a wavelength dispersion modifier for polymer film (i.e., JP-A-2009-64006 and JP-A-2009-64007). JP-A-2009-270062, JP-A-2009-79213 and JP-A-2009-67983 have proposed techniques of using a merocyanine-based ultraviolet absorber combined with other types of ultraviolet absorber, aiming at improving the light stability.\nPolymer film typically used as an optically compensatory film of liquid crystal display device is sometimes required to show high Rth. Increase in the amount of addition of the additive, aimed at elevating Rth, have however occasionally degraded the stability of manufacturing, due to deposition of the additive onto the surface of the film in the process of film making. In the field of manufacturing of polymer film typically adoptable to optical films, it is therefore very beneficial to provide an additive capable of ensuring high Rth even with a less amount of addition, while ensuring a desired level of wavelength dispersion."}
-{"text": "An 802.16 family of standards have been developed by the Institute of Electrical and Electronic Engineers (IEEE) to provide for fixed, portable, and/or mobile broadband wireless access (BWA) networks (for example, the IEEE std. 802.16, published 2004 and subsequent revisions). The Worldwide Interoperability for Microwave Access (WiMAX) forum facilitates the deployment of broadband wireless networks based on the IEEE 802.16 standard. In particular, the WiMAX forum ensures the compatibility and inter-operability of broadband wireless equipment.\nIn downlink transmissions, WiMAX networks may broadcast data packets from a base station (BS) to a subscriber station (SS) or mobile station (MS), whereas in the uplink transmissions, the scheduling services may be designed to support services with different traffic characteristics and Quality of Service (QoS) requirements. A significant benefit of the converged wireless networks, such as a WiMAX network, is in the sharing of the most valuable resources\u2014the wireless spectrum among different services. However, the wireless network convergence in a WiMAX network also comes with some challenges, due to the arbitration of uplink transmission between multiple SSs, as well as the allocation of uplink bandwidth with QoS needed for different services.\nIt is desirable to have a system and method for adaptively updating wireless uplink scheduling based on characteristics of the uplink data, to improve a user's experience, while minimizing impact on the wireless spectrum."}
-{"text": "Automatic card shuffling machines were first introduced by casinos approximately ten years ago. Since then, the machines have, for all intents and purposes, replaced manual card shuffling. To date, most automatic shuffling machines have been adapted to shuffle one or more decks of standard playing cards for use in the game of blackjack. However, as the popularity of legalized gambling has increased, so too has the demand for new table games utilizing standard playing cards. As a result, automatic shuffling machines have been designed to now automatically \u201cdeal\u201d hands of cards once the cards have been sufficiently rearranged.\nFor example, U.S. Pat. No. 5,275,411 (\u201cthe '411 patent\u201d) to Breeding and assigned to Shuffle Master, Inc., describes an automatic shuffling and dealing machine. The '411 patent describes an automatic method of interleaving cards as traditionally done in a manual fashion. Once interleaved, the entire stack of shuffled cards is positioned above a roller that removes and expels a predetermined number of cards from the bottom of the stack to a card shoe. Once the predetermined number of expelled cards are removed from the shoe by a dealer, a second set of cards is removed and expelled. This is repeated until the dealer has dealt each player his or her cards and has instructed (e.g., pressed a button on the shuffler) the shuffling machine to expel the remaining cards of the stack.\nThe '411 patent and related shufflers, having a dealing means, suffer from the same shortcomings\u2014slowness, misdeals and failure. However, the machines currently marketed are still favored over manual card shuffling. On the other hand, since casino revenue is directly proportional to the number of plays of each wagering game on its floor, casinos desire and, in fact, demand that automatic card shufflers work quickly, reliably and efficiently.\nAccordingly, the present invention utilizes a proprietary random card ejection technique, in combination with a novel card separation and delivery unit, to overcome the aforementioned shortcomings. The present invention uses random ejection technology to dispense individual cards from a card input unit to a card separation and delivery unit of the shuffler. A card stop arm and floating gate control the number of ejected cards that may, at any one time, travel to the card separation and delivery unit. The ejected cards are then separated by a feed roller system which propels the cards to a collection unit. Once a predetermined number of cards are propelled to the collection unit, additional cards are ejected from the card input unit. A shuffler processing unit in communication with internal sensors controls the operation of the shuffler.\nAn audio system is adapted to communicate internal shuffler problems and shuffler instructions to an operator. Preferably, the audio system is controlled by the shuffler processing unit in communication with a second local processing unit."}
-{"text": "In contrast to semiconductor production of silicon wafers, which are processed individually, the manufacture of devices such as micromechanical sensors is accomplished using relatively small substrates. In this context, the base substrate is present in the form of a small cylindrical steel element. In order for the manufacturing process and for application of a circuit produced in thin-film technology onto such substrates to be cost-effective, the substrates are processed in multiple fashion. The individual base substrates are therefore retained in large numbers on a workpiece carrier, and processed together.\nHigh-pressure sensors that are utilized in many automotive engineering systems and in automation technology are fabricated on such substrates. Applications include, for example, direct fuel injection, common-rail technology for diesel vehicles, electrohydraulic brake systems, vehicle dynamics control systems, and many more.\nPressure is sensed via the deflection of the sensor membrane, which is coated with a Wheatstone measurement bridge using thin-film technology. The circuit is patterned using a photolithographic manufacturing process. Contacting and passivation can also be performed using a production step of this kind. The quality of high-pressure sensors depends substantially on the accuracy with which the resistance features of the Wheatstone measurement bridge are positioned on the pressure membrane. The resistances must be centered as much as possible on the pressure membrane. This positional accuracy substantially influences the electrical properties of the measurement bridge, for example, the range of the signal furnished by the sensor.\nIt is believed that the photolithographic process of patterning the functional layer to form the resistances is therefore essential in the production of a high-pressure sensor. A preconditioning of the surface of the functional layer, application of a photosensitive resist onto the functional layer, an oven step to condition the resist for exposure, exposure of the resist through an exposure mask, development of the exposed features, an oven step to condition the developed resist structure for etching, mapping of the resist image into the functional layer using an etching process, and subsequent stripping of the photoresist mask from the completed etched functional structure, are usually performed in this context.\nThe location of the resistances may be determined in the exposure step. The substrates, previously diced and retained on a workpiece carrier, may then be exposed with sufficient positional accuracy relative to an exposure mask. In contrast to the exposure of silicon wafers in semiconductor production, in which the wafers can be registered as one piece with the exposure mask (e.g. a quartz mask) to be imaged, in the exposure of the previously diced plurality of substrates there may be a positional inaccuracy for each substrate with respect to the imaging exposure mask.\nIt is believed that the substrates are often handled in immovably bolted-down workpiece carriers that must each be individually registered with respect to the exposure mask. This registration operation is performed manually.\nFor example, in the projection exposure method of the Nagano Keiki Co. Ltd., a special workpiece carrier that is automatically positionable is provided in each case. This workpiece carrier is not identical for the subsequent processes."}
-{"text": "With the development of intelligent mobile terminals and the popularity of social communication networks, the social communication scope of an intelligent mobile terminal user is greatly expanded. Therefore, the mobile terminal user has lots of contact information to manage. The current contact management manner has some obvious disadvantages in managing large amounts of contact information. For example, a contact is difficult to search for or find because the contact is generally searched for by using the first letter of an alphabetic identifier of the contact, that is, inputting the first letter of the contact in a contact search box. Because the contact information is huge, a large number of contacts having the same first letter may be searched out or found by using the first letter. These contacts having the same first letter are all displayed. In this case, the user is only capable of finding needed contact information by browsing or scrolling pages. This process is time-consuming, and does not satisfy a user's expectation for quickly searching for and locating a contact. In addition, user experience on the intelligent mobile terminal is reduced."}
-{"text": "Consumption of nutrients, like antioxidants and folic acid, which are abundant in fruits and vegetables, has been linked to a lower incidence of cardiovascular disease. Moreover, it is well settled that eating fruits high in soluble fiber can reduce cholesterol levels which protects against atherosclerosis.\nOther advantages of having a diet high in fruit include better athletic performances, reduced risk of developing chronic bronchitis, a lowered risk of getting most common cancers (including breast cancer), as well as a lowered risk of getting cataracts. Additionally, fruits, like avocado, are high in protein, vitamin E, C, and beta-carotene and may be linked to stroke prevention.\nWhile food products comprising fruits and vegetables have been linked to health benefits in humans, such products are often difficult to prepare for sale in commerce. This is true because the quality of food products comprising fruit often deteriorates (e.g., browns, darkens, grows mold and/or loses flavor) due to enzymatic reactions within the food product, thereby resulting in a product that has a short shelf life and does not have an appealing look or taste after spending a limited period of time in conventional commercial channels.\nKnown techniques have been used to inhibit the deterioration of food products comprising fruits and/or vegetables. These known techniques include pasteurization of the fruit, high vacuum processing for removing oxygen, and chemically treating the fruit with sulfiting agents before making the food product. The above-described known techniques do not eliminate, for example, browning and darkening in food products comprising fruit, and such techniques have adverse effects on the flavor, aroma, texture and nutritional value of the fruits and vegetables treated, as well as the food products prepared therefrom.\nIt is of increasing interest to develop a ready-to-serve and ambient stable fruit-based composition that does not, for example, easily brown, darken and lose flavor and that has an extended shelf life at ambient temperature. This invention, therefore, is directed to a ready-to-serve and ambient stable fruit-based composition that has not been subjected to chemical treatment, high vacuum processing and temperatures over about 85\u00b0 C. The ready-to-serve ambient stable fruit-based composition of this invention is low in fat, substantially free of starch, and suitable to be a beverage, filling, dip, sauce, spread, dressing or the like. Moreover, the ready-to-serve ambient stable fruit-based composition of this invention can be prepared from under ripe, ripened and/or over-ripened fruit and is ambient stable for at least about eight (8) weeks after opening. Such a fruit-based composition has substantially the same visual, texture, aroma and taste attributes of a fruit-based composition made on demand from freshly picked fruits."}
-{"text": "The present invention relates generally to porous materials having high surface areas, and more particularly to such materials which are non-degradable, and methods for fabricating the same.\nMembranes have been typically used for filtration (microfiltration, ultrafiltration, nanofiltration), reverse osmosis (hyperfiltration), dialysis, pervaporation, and gas separation applications. See Scott, K. and R. Hughes, Industrial Membrane Separation Technology, 1996, London: Blackie Academic & Professional; Baker, R. W., Membrane technology and applications, McGraw-Hill professional engineering, 2000, New York: McGraw-Hill; and Cardew, P. T., M. S. Le, and Royal Society of Chemistry, Process Technology Group, Membrane processes: a technology guide, 1998, Cambridge: Royal Society of Chemistry.\nA solid membrane can be made of synthetic polymers, natural macromolecules, inorganic compounds, ceramic, or metallic materials. These membrane materials are generally fabricated through sintering, stretching, extrusion, phase inversion and etching, or casting. See Scott, K. and R. Hughes, Industrial Membrane Separation Technology, 1996, London: Blackie Academic & Professional; and Pinnau, I. and B. D. Freeman, Membrane formation and modification, ACS symposium series, 744, 2000, Washington, D.C.: American Chemical Society, Distributed by Oxford University Press.\nPorous membranes are advantageous in their low resistance to mass transfer of solutes in solution due to the increased permeation rate resulting from the pores. Therefore, porous membranes have been employed for separation of mixtures of proteins and macromolecules, salt concentration, sterilization, etc. They can also serve as 3-D matrices for chemical and biochemical mass exchange or reactions to take place, or for cells or other living organisms (e.g., bacteria, viruses, fungi) to grow. Therefore, they can be used as matrices in diagnostic systems, catalysis systems, culture systems, drug delivery systems, wound dressings, etc.\nU.S. Pat. No. 6,146,892 discloses a method for producing nanofibrillar matrices utilizing degradable polymers, such as for example, poly(L-lactic acid) (PLLA), poly(D,L-lactic acid-co-glycolic acid) (PLGA), and the like. The disclosed nanofibrillar matrices are highly porous and work well for various applications. However, the pore sizes are very small (on the order of 2 \u03bcm to 3 \u03bcm), which may in some instances render it more difficult for cells to enter. Further, small pore sizes may render it more difficult for material transport, especially materials which are particulate or contain particles. Yet further, the disclosed non-fibrillar structure (a platelet structure) was not as mechanically strong as may be desirable in some instances."}
-{"text": "1. Field of the Invention\nThe present invention relates to a document data management system.\n2. Description of the Related Art\nIn an environment in which a document management system is used, a data of a document which is received via facsimile, a document which is scanned and transmitted by a multifunction peripheral or the like can be managed by storing the data to a folder of the document management system. A user carries out various kinds of operations with regard to the stored document data by making an access to a document management server from a client PC to display the prescribed folder. In this case, if a user carries out an operation which is the same as that he/she carried out in the past with regard to, for example, a stored facsimile document data, the user needs to select the operation again and execute the operation. That is, a user needs to repeat the same operation.\nHence, there has been proposed a document management system which can record information of an operation with regard to a document data as related information of a document, and select and acquire the information when other document data is operated and execute the operation.\nThere is a case in which a user combines a document data which is stored in a device of a multifunction peripheral or the like, and other document data which is stored in a document management system to be a single document, and carries out a collective operation with regard thereto. As a collective operation, for example, plural document data (file) are combined and converted into PDF, plural document data are made to be a binder, plural document data are summarized and transmitted by facsimile or an electronic mail and so on. However, in a case in which document data which are stored in different folders are summarized and subjected to a collective operation, it is necessary to copy the respective document data temporarily in a work folder (temporary work area), which takes time and labor. Also, in a case of carrying out the collective operation as described above with regard to a document data stored by using a sort server, it is conceivable that a combination of document data which are subjected to the collective operation differs for respective individual document data. In this case, it is necessary to carry out the operation by collecting document data of an object of the collective operation in accordance with respective document data.\nFor example, according to a technology which is described in Japanese Patent Laid-Open No. 2009-230221, although related information can be applied to a single document data, it is necessary to acquire and select the related information and execute the application at respective times thereof, and therefore, time and labor are taken."}
-{"text": "The present invention relates to a master cylinder configured to supply a liquid pressure to a braking cylinder of a vehicle.\nThis application is the U.S. national phase of International Application No. PCT/JP2013/058948 filed 27 Mar. 2013 which designated the U.S. and claims priority to Japanese Patent Application No. 2012-241198, filed Oct. 31, 2012, the contents of each of which are incorporated herein by reference."}
-{"text": "A sheet of the type of the present invention is disclosed by German Auslegeschrift 40 18 716. It is manufactured using a rubber-elastic mixture of special Shore A hardness of less than 85 which can be processed thermoplastically, and contains a thermoplastic polyurethane elastomer (TPU) and an ethylene/vinyl acetate copolymer (EVA). The vinyl acetate (VA) content of the EVA is taught to be between 78 and 95% by weight. In the rubber-elastic mixture there should be 50 to 99% by weight of TPU and 1 to 50% by weight of EVA, wherein the sum of the percentages is equal to 100%. The mixture is used to produce molded objects and sheets which are soft and elastic, and are readily processed (that is, without sticking), yet they also have a high strength, a good elongation at break, a high abrasion resistance, and show little swelling in fuels and lubricants.\nThe VA content of the EVA materials which are presently commercially available is of the order of 40 percent by weight. These materials can be calendered. If such an EVA is used in a thermally-processible mixture in accordance with the information of the German Auslegeschrift 40 18 716, the resulting mixture cannot readily be processed into sheets by calendering because the melt strength is too low. Moreover, calendering of this mixture results in a sheet with unacceptable sticking qualities.\nIt is therefore an object of the invention to develop an improved thermoplastic material having the desirable properties of the material of German Auslegeschrift 40 18 716, as well as having the capability of being processed into sheets, for example, by calendering."}
-{"text": "The activities of a business or enterprise can be grouped into processes. Processes are business operations that are separated as desired and usually occur across business units. For example, the process of taking orders and turning those orders into revenue may be known as Order to Cash. The processes are comprised of sub-processes. For example, Order to Cash may be broken down into sub-processes such as Receive Order Inquiry, Provide Customer Quotation, Create Customer Outline Agreement, Create Sales Order, Schedule Production, Manufacture Product, Ship Product and Invoice Customer. Each sub-process may in turn be broken down into discrete activities such as providing customer number, entering that customer number, establishing pricing, determining a shipping date, etc.\nThe processes, sub-processes and activities operate, in part, by communicating information. For example, users may communicate through email. As another example, applications may communicate amongst themselves through electronic data interchange (\u201cEDI\u201d) and other similar services. Communication occurs horizontally, that is, among a process, sub-process and activities, as well as vertically, that is, between processes, sub-processes and activities.\nWhether communications occur horizontally or vertically, among applications or users, communications are increasingly asynchronous or message based. That is, enterprise communications were formerly primarily synchronous, or connection oriented, in which a connection is established with prior coordination between communication end points with data then being transmitted over the connection. Enterprise communications are now increasingly asynchronous, or connectionless, transmitting data without prior coordination between communication end points, such as through \u201cevent based\u201d communications which use messages to move data instead of large files.\nAsynchronous or message based communications permit loosely coupled connections among and between systems because the end points do not have to be prepared to receive the data when the message is transmitted. Loosely coupled connections permit more flexibility in assembling processes. Flexibility in assembling processes is desirable in order to permit quick reaction to changing business conditions: if a particular sub-process or activity becomes unusable, the process can be reassembled with a new sub-process or activity. For example, if a Manufacture Product sub-process in the Order to Cash process at Widget Co. enterprise has a specific factory identified to manufacture the product and that factory has a fire or other disaster, making it unusable, Widget Co. will need to substitute a new factory. The ripple effect of that substitution among all of Widget Co.'s processes will change any number of parameters. A loosely coupled asynchronous connection among Widget Co.'s processes provides rapid substitution of the new factory for the old because the end points of communication to the new factory do not have to be predetermined before communications begin with the new factory. Thus, the flexibility of the asynchronous message based communication has permitted quick response to changing business conditions.\nDespite this flexibility, asynchronous or message based communications are problematic because of their loosely coupled nature. At any given time, precise information on the progress of the processes is difficult to obtain\u2014messages may be in transit and not instantly locatable. For example, if a customer calls for the status of an order, an enterprise customer service representative may be able to determine nothing more than the fact that the order has been received and that the scheduled ship date is X. There is often no ability to drill down into the information levels and review the status in more granularity, such as the location of the good, the manufacturing status, etc., because the information required to review that status is in transit and unable to be reviewed.\nOf course, if the enterprise lacks the ability to access status information, business partners of the enterprise will similarly lack that ability. Thus, asynchronous communications may well increase inefficiency among business partners as well.\nThe difficulty in reporting caused by message based architecture also makes it difficult for the enterprise to measure the efficiency of its processes and their sub-process. Asynchronous messaging, with its indeterminate transmission of information, means a company may not be able to easily measure the interval between each sub-process, e.g. the time between Scheduling Production and the Manufacturing of a Product, and so easily measure the efficiency of their operations.\nFinally, asynchronous messaging may provide an enterprise with an ability to model and simulate processes. That is, since information flows can be readily estimated through enterprises with asynchronous messaging, and processes can be easily modeled from those flows, asynchronous messaging modeling provides the potential to model and simulate processes. That potential is not realized with present technology, however. Moreover, since as described above, enterprises lack information on the processes they have implemented, the enterprises are handicapped in their ability to modify those processes or plan new processes. A modeling and simulation tool, demonstrating processes, sub-processes and their activity or granular detail level would be extremely helpful, and would give the enterprise an opportunity to assemble, test, adjust, and simulate processes and their details.\nAccordingly, it is an object of the present invention to provide a tool for simulating message based architectures.\nIt is a further object of the present invention to provide monitoring capabilities for enterprise processes."}
-{"text": "The present disclosure relates generally to a downhole nuclear magnetic resonance (NMR) apparatus, data processing, and interpretation methods for evaluating a characteristic of a region, and particularly for detecting and quantifying a gas-bearing earth formation in a subterranean region.\nNMR well logging is a technique used to investigate subterranean regions that may contain water, oil and/or gas reserves. The nuclei of chemical elements have a characteristic angular momentum (spin) and a magnetic moment, and by detecting and analyzing the reaction of the nuclei to applied magnetic fields, the characteristics of specific nuclei may be deduced. In the presence of an externally applied static magnetic field (B0), the nuclei spins become magnetized and align themselves parallel to the B0 field. By applying a radio frequency (RF) pulse train of a specific frequency to the magnetized nuclei, a pulsed RF magnetic field (B1) is generated that tips, or flips, the spins away from the direction of the B0 field. If the RF frequency (\u03c9) substantially matches the condition for NMR (\u03c9=\u03b3B0), where \u03b3 is the gyromagnetic ratio, then the first pulse (herein referred to as A pulse) reorients the magnetization to start precession and subsequent pulses (herein referred to as B pulses) generate spin-echo signals. A RF pulse sequence known as the CPMG (Carr-Purcell-Meiboom-Gill) sequence is typically used for well logging.\nAt the end of an A pulse, the spins are oriented transverse to the B0 field and precess around the direction of the B0 field at the Larmor frequency (\u03c90=\u03b3B0), and the transverse magnetization dephases with a transverse relaxation time constant (T2), also known as the spin-spin relaxation time. Repeated tipping and relaxation of the spins results in the NMR spin-echo signal, which may then be detected and analyzed for oil and/or gas field exploration.\nExisting methods use dual wait-time logs and the T1 contrast between gas and other formation fluids for gas detection and for gas saturation estimation. One such method is based on the assumption that water signals are fully polarized at both short and long wait time, TW, but gas signal is only partially polarized. So the difference between the two is contributed from gas only. However, if three phases coexist in the NMR detected sensitive volume, especially if slow relaxing water signal and light oil or oil-based mud filtrates are present in the formation, detection may be difficult or limited. Other methods for acquiring and processing multiple wait time data for T1 estimation employ improved log quality data using a summation of echoes approach, which is more useful if all of the partially polarized signal is gas. However, the certainty of the signal is reduced for discerning gas using this technique for wells drilled with OBM (oil based mud), and the subsequent invasion of the OBMF (oil based mud filtrate) into the sensitive volume, or for formations containing large-pore water or light oil. Thus skilled human interpretation is required in order to use existing art techniques. Accordingly, there is a need in the art for a robust NMR detection and analysis method that overcomes these drawbacks."}
-{"text": "As a result of progress in current CMOS (complementary metal oxide semiconductor) technology, such as advances in 65-nm and 45-nm processes, in a programmable logic device (PLD), such as a field programmable gate arrays (FPGAs), have increased logic elements and I/O buffer (or pin) densities. Because simultaneous switching noise (SSN) is associated with the number of toggling I/O buffers, signal voltage level, and the switching rate of I/O buffers, it is important for FPGA users to be able to evaluate SSN performance in a \u201cchip-package-board\u201d environment. The SSN caused by FPGA output buffers is widely known as simultaneous switching output (SSO) noise and can be differentiated from the SSN caused by input buffers.\nSSN in an FPGA system may be attributed to two primary factors: the mutual inductive coupling among switching I/O buffers and the impedance profile of a power distribution network (PDN) including die, package, and printed circuit board (PCB). In essence, reducing SSN is a design-cost issue. While designers can minimize the mutual inductive coupling by increasing the ratio of FPGA ground pins (or return-current pins) to I/O buffers, this approach sacrifices I/O buffer densities. Engineers can improve PDN performance by increasing on-die capacitance and adding on-package decoupling capacitors, but this approach increases costs.\nBecause FPGAs are programmable, they fit into a wide variety of user applications, so it is useful for designers to have a tool to determine their own SSN budget without additional costs. Ideally, this kind of tool requires instantaneous and accurate result predictions for various FPGA I/O buffer assignments and ideally would provide an optimum buffer assignment under certain design constraints. In the past, designers have constructed system-level, \u201cSPICE-like\u201d models to anticipate SSN in FPGA systems. These models are based on an understanding of SSN cause mechanisms and correlate well with bench measurements, helping IC and packaging designers improve designs.\nHowever, these models are so complicated that they require signal/power integrity expertise to perform time-consuming, system-level simulations. It is cumbersome for all FPGA designers to perform the same level SSN analysis without considering their different design margins. Therefore, there is a need for a tool to help FPGA users execute a comprehensive FPGA SSN analysis in a short design cycle, though the invention claimed below has applicability to other applications beyond this particular application, as will become apparent from the following description and the drawings which accompany it."}
-{"text": "Lately with the development of the housing industry, the secondary housing industry has progressed very rapidly, including manufacture of built-in furniture, kitchen sinks, washing stands as unit equipment.\nKitchen sinks and washing stands are one of the interior fixtures essential for every day life, and a variety of such products directed to the taste of customers are being presented on the market.\nSince kitchen sinks and washing stands are fixtures which inevitably involve use of water, a top which supports a basin is provided horizontally as one of the constituent members. As such a top there is usually employed that having a structure obtained by making the surface of a wood derived core material repellent to water and at the same time giving a color so as to serve also as surface decoration or bonding a melamine decorative board having a woodgrain pattern. The kitchen sinks and washing stands are, as described above, fixtures utilizing water indoors, and therefore the used water should not leak into other areas than the sinks and washing stands.\nAccordingly, various devices have been imparted to the aforementioned tops fitted with basins in order to prevent leakage of water spilt on the surface of the top through a joint part between a flange part of the basin and the edge of a hole in the top for imbedding the basin. One such example is a product in which a basin flange part is mounted on a top, the joint part thereof with the edge of a hole in the top for imbedding the basin is sealed by fixing a rubber packing or other sealing materials, and the top and the basin flange part are fixed together by screwing from the rear surface of the top. However, those having such a structure suffer from disadvantages that with prolonged use, the rubber packing or other sealing materials are degraded and allow spaces around the edge of the aforementioned imbedding hole to cause water leakage, thereby resulting in swelling and corrosion of the top itself, which then forms gaps between the surface of the top and the basin flange part, finally permitting water spilt on the top to drop onto the floor instead of flowing into the basin, thus making the neighborhood of the sink or washing stand very moist.\nAs an alternative to those having such structure that the joint part between the basin flange part and the edge of the hole in the top for imbedding a basin is sealed using a rubber packing or other sealing materials, there is a structure in which a hole in the top is made larger than the basin itself, and the top decorative board alone is provided with a hole having the same dimention as the open part of the basin, and the basin flange part is fitted and directly adhered to the hole edge of the decorative board on its rear surface using a synthetic resin. By being of such a structure, since the flange part of the basin is situated lower than the surface of the top, all the water spilt on the surface of the top flows into the basin, thus solving the aforementioned problem of drainage of water spilt on the surface of the top. However, with this structure, the strength of the bonded part between the top decorative board and the basin flange part is weak, and it therefore tends to be damaged during transport and the decorative board is also easily impaired. Moreover, the finish treatment of the synthetic resin used for adhering the joint part between the hole edge of the decorative board and the basin flange part is troublesome and requires a lot of labor and time, and therefore it is not applicable to manufacture of products having a structure suitable for mass production.\nAs mentioned above, in fabricating tops fitted with basins used in kitchen sinks, washing stands etc., various problems are present from a standpoint of working to prevent water leakage in the joint part between the basin flange part and the hole in the top for imbedding a basin.\nThe present invention aims to present tops fitted with basins which can provide good water drainage and have a structure permitting their mass production by eliminating problems involved in the structure and fabrication of tops fitted with basins which are constituent members for kitchen sinks and washing stands, as well as a process for their production."}
-{"text": "With the wide application of Chinese characters, learners of Chinese often fail to pronounce some word segments in a process of reading a Chinese character text aloud, especially some polyphonic word segments. Therefore, the learners of Chinese always come across obstacles of reading aloud. Even though pronunciations of a polyphonic word segment can be obtained by looking up a dictionary, much time is required to do so. In addition, because a vocabulary included in the dictionary is large, it is difficult for a user to grasp the most accurate pronunciation of the polyphonic word segment. This lowers the efficiency and effects in text reading, thereby affecting effects in Chinese learning."}
-{"text": "In 1962, Hanna-Barbera Production, Inc. produced a cartoon titled \u201cThe Jetsons\u201d that featured the Jetson family living in a utopian future having incredible conveniences based on communication devices, large and small displays, and even robots. Everything they wanted required only a voice command or a push of a button.\nIn 1966, a television show titled \u201cStar Trek\u201d debuted that followed the interstellar adventures of James T. Kirk and the crew of the Enterprise. Star Trek is known for its influence on the world outside of science fiction for inspiring inventions such as desktop and tablet computers, wireless head phones, biometrics, cell phones, and even the automated sliding door.\nAn epic science fiction film released in 1968 titled \u201c2001: A Space Odyssey\u201d follows a voyage of a U.S. spacecraft to Jupiter. The film had tablet computers like Star Trek but also had two-way video conferencing, a suitcase phone, a voice controlled chessboard, and a sentient computer named HAL having artificial intelligence.\nForty-three years later, Apple\u00ae introduced the iPhone 4S\u00ae, which includes a computer program called Siri\u00ae that works as an intelligent personal assistant and knowledge navigator. Siri uses a natural language user interface to answer questions, make recommendations, and perform actions. Siri adapts to a user's individual language usage and preferences and returns individualized results.\nMany vehicles now offer various features activated by voice commands such as making telephone calls, selecting music, and interfacing with car navigation systems to find nearby restaurants, hotels or gas stations or to receive driving directions to a known address. Pioneer sells a car receiver that connects to an iPhone to take advantage of Siri\u00ae Eyes Free functionality including listening to text messages, calendar information, and reminders.\nA home automation system integrates electrical devices in a house with each other. The techniques employed in home automation include those in building automation as well as the control of domestic activities, such as home entertainment systems, houseplant and yard watering, pet feeding, changing the ambiance \u201cscenes\u201d for different events (such as dinners or parties), lighting control system, and the use of domestic robots. Devices may be connected through a home network to allow control by a personal computer, and may allow remote access from the internet. Through the integration of information technologies (e.g., PCs, phones) with the home environment, systems and appliances can communicate in an integrated manner which results in convenience, energy efficiency, and safety benefits.\nAutomated \u201chomes of the future\u201d have been staple exhibits for World's Fairs and popular backgrounds in science fiction. However, problems with complexity, competition between vendors, multiple incompatible standards, and the resulting expense have limited the penetration of home automation to homes of the wealthy or ambitious hobbyists.\nCurrent voice-controlled home automation systems typically involve wall-mounted control panels that provide a combination of touch and voice-controlled features. Such systems typically require a person to be standing next to the control panel, which substantially limits the \u2018value add\u2019 of using voice commands over using touch features of the control panel or a phone.\nTechnology improvements are needed to remove current limitations and inefficiencies of home (and business) automation systems."}
-{"text": "Much research has been carried out in the field of human probiotics in the last decade (see review Huis in't Veld et al. (1994) Tibtech 12, 6\u20138). This research has been prompted by the rising interest by the public in their health and well-being. Many probiotic products are now available on the market and some of the beneficial effects derived from these products range from alleviation of lactose intolerance (Gilliland, S. E. (1990) FEMS Microbiol. Rev. 87, 175\u2013188) to prevention of diarrheal diseases (Marteau, P. et al. (1993) FEMS Microbiol. Rev. 12, 207\u2013220) and possible prevention of carcinogenesis (Adachi, S. (1992) In \u201cThe Lactic Acid Bacteria in Health and Disease\u201d. (Wood, Ed.), 233\u2013262, Elsevier, Barking). Controversy exists over many of these beneficial effects as no standardised procedures are available and contradictory results have been published with regard to the possible beneficial effects of cultured products containing \u2018probiotic\u2019 bacteria.\nPoor choice of strain has been cited as one of the contributing factors to the inconsistency and variability of results (Marteau, P. et al. (1993) supra) (Kim, H. S. (1988) Cult. Dairy Prod. J. 23, 6\u20139) and Fuller, R. ((1989) J. Appl. Bact. 66, 365\u2013378) outlined criteria pertaining to the successful isolation of probiotic strains. The strains should be indigenous to the intended host species and also have the ability to (i) survive and grow within that host; (ii) exert a beneficial effect at the target site and (iii) be maintainable in the carrier food or system throughout product manufacture and storage.\nThere is a fast growing market for health-promoting products including probiotics. Many such products are now available (Jong, S. C. and Birmingham, J. M., (1993) ATCC Quart. Newslett. 13(1), 1\u201311). One of the more important components of these products is the microorganisms used. The most frequently utilised species include Bifidobacterium sp., Lactobacillus sp., and Propionibacterium sp. (O'Sullivan, M. G., et al. (1992) Trends in Food Sci. and Tech. 3(12), 309\u2013314). There is a lack of substantiated evidence from controlled trials that the organisms currently used in such products are those which have beneficial effects on the gut flora (Tannock, G. W. (1983) In Human Intestinal Microflora in Health and Disease 517\u20135399 D. J. Hentges (ed.), New York, Academic Press). The source of the microorganism is critical to its survival and therefore its function in the human intestinal tract. Lee, Y-K and Salminen, S. ((1995) Trends Food Sci. Technol. 6, 241\u2013245) stated that as a general requirement, a probiotic strain should be of human origin as some health-promoting effects may be species dependent. It is well known that the indigenous microflora is one of the major defense mechanisms that protects the human against colonisation by allochthonous invading bacteria (Tancrede, C. (1992) Eur. J. Clin. Microbiol. Infect. Dis. 11(11), 1012\u20131015) and it is also the human's best ally when supporting the immune system. Bacterial populations at different levels of the gastrointestinal tract constitute complex ecosystems depending on the physiology of the host and on interactions between bacteria.\nTen Brink et al. ((1994) Journal of Applied Bacteriology 77 140\u2013148) isolated and screened a large number (\u02dc1000) of Lactobacillus strains for the production of antimicrobial activity. Lactobacilli were isolated from various fermented foods and feeds (sauerkraut, cheese, sausage and silage), human dental plaque and faeces derived from different laboratory animals (rat, mouse, guinea pig and quail) and human volunteers. Only eight positive strains were found and two of these were studied, namely Lactobacillus salivarius M7 and Lactobacillus acidophilus M46. The former strain produces the broad spectrum bacteriocin salivaricin B which inhibits the growth of Listeria monocytogenes, Bacillus cereus, Brochothrix thermosphacta, Enterococcus faecalis and many lactobacilli. L. acidophilus M46 produces a bacteriocin acidocin B which combines the inhibition of Clostridium sporogenes with a very narrow activity spectrum within the genus Lactobacillus. However, these strains are not indigenous to the infected host species, which is one of the criteria which is required for a successful probiotic strain for human use.\nArihara, K. et al. ((1996) Letters in Applied Microbiology 22, 420\u2013424) have isolated Salivacin 140 a bacteriocin from Lactobacillus salivarius subsp. salicinius T140. Strain T140 was isolated from the surface of Japanese pampas grass leaves grown close to an animal barn and thus the strain was likely to have derived from animal faeces.\nThere is a need for probiotic strains which meet the aforementioned criteria. Bacteriocin production by lactobacilli is thought to play an important role in the competitive exclusion of pathogens and other undesirable microorganisms of the intestinal tract of humans. Bacteriocins are broadly defined as proteinaceous compounds which exhibit a bactericidal effect against a wide range of microorganisms.\nDue to their diversity of species and habitats lactobacilli are the most bacteriocinogenic of the lactic acid bacteria. As many as forty bacteriocins produced by lactobacilli have now been isolated (Klaenhammer, T. R. (1993) FEMS Microbiol. Rev. 12, 39\u201386).\nBacteriocins have been isolated from human infant faeces. However, the bacteriocins were found to have narrow host ranges and were active only against other lactobacillus species (Toba, T. et al. (1991) Lett. Appl. Microbiol. 12, 228\u2013231.).\nThere is a need for bacteriocins with a broad spectrum of activity."}
-{"text": "The current disclosure relates generally to foundations for mattresses to increase mattress support, and more particularly to foundations having novel slat construction to improve mattress support.\nWhile traditional foundations are sufficient in many situations, applicant believes that in some situations, improvement can be made. By way of example, in some situations, a sleeper's downward force on traditional slats may cause the mattress to dip in between the slats of the foundation (referred to herein generally as \u201cmattress dip\u201d). Applicant believes mattress dip may be caused by any of a number of problems, including at least one of mattress spring position relative to the slat, spring size, slat size, slat spacing, sleeper weight, sleeper habits, etc. Regardless of the cause, applicant believes mattress dip can result in decreased sleeping comfort and/or mattress life.\nAccordingly, applicant desires foundations that can reduce mattress dip by providing increased mattress support."}
-{"text": "The desired characteristics for electrode materials which are to be used in redox cells to and other electrolytic processes include: low electrical resistivity, good mechanical properties and chemical stability and optionally low weight and low volume. The electrode material should not be attacked by reactants or products during cell operation. Thus there is a need for an electrode material which has good mechanical properties to withstand the hydraulic pressure and cell assembling operation without cracking and which is solution impermeable so that penetration and attack of the metal current collector (in end-electrodes) and loss in current efficiency (due to electrolytes crossing the bipolar electrodes) is low."}
-{"text": "As an exemplary non-volatile semiconductor memory device, a type of non-volatile semiconductor memory device has been conventionally known that accumulates electric charges into a charge accumulation layer of a memory cell transistor by means of, for instance, a quantum tunneling effect in order to execute data writing (e.g., see PTL 1). Actually, as shown in FIG. 6, a non-volatile semiconductor memory device 100 of this type has a structure that higher order bit lines 101a and 101b and word lines 102a to 102h are disposed in an intersecting manner and a plurality of memory cell transistors 103 are disposed in a row and column matrix with respect to the higher order bit lines 101a and 101b and the word lines 102a to 102h. \nThe higher order bit line 101a is provided with a plurality of first semiconductor switches 104a and 104c, while a single lower order bit line 105a, 105c is connected to each first semiconductor switch 104a, 104c. Further, in this exemplary embodiment, the other higher order bit line 101b is also similarly provided with a plurality of first semiconductor switches 104b and 104d, while a single lower order bit line 105b, 105d is connected to each first semiconductor switch 104b, 104d. In such non-volatile semiconductor memory device 100, each lower order bit line 105a, 105b, 105c, 105d forms a memory block 106a, 106b, 106c, 106d, while each memory block 106a, 106b, 106c, 106d has the plural memory cell transistors 103.\nThe first semiconductor switches 104a, 104b, 104c and 104d are herein formed by N-MOS (Metal-Oxide-Semiconductor) transistors. Further, the first semiconductor switch 104a of the memory block 106a, for instance, is connected at the source thereof to the higher order bit line 101a while being connected at the drain thereof to the lower order bit line 105a, and is also connected at the gate thereof to a first selected gate line 108a shared with another memory block 106b aligned with the memory block 106a along a row direction. Thus, by this configuration, a predetermined gate voltage can be equally applied from the shared, single first selected gate line 108a to the two first semiconductor switches 104a and 104b that are mounted on the memory blocks 106a and 106b disposed in the upper part of FIG. 6.\nOn the other hand, similarly also in the two memory blocks 106c and 106d aligned along the row direction in the lower part of FIG. 6, a single first selected gate line 108b is connected to the two first semiconductor switches 104c and 104d, and a predetermined gate voltage is can be configured to be equally applied to the two first semiconductor switches 104c and 104d from the shared, first selected gate line 108b. \nIn addition to this, a higher order source line 110a is provided with a plurality of second semiconductor switches 111a and 111b, while a single lower order source line 112a, 112b is connected to each second semiconductor switch 111a, 111b. Further, the other higher order source line 110b is also similarly provided with a plurality of second semiconductor switches 111c and 111d, while a single lower order line 112c, 112d is connected to each second semiconductor switch 111c, 111d. \nFurther, the second semiconductor switches 111a, 111b, 111c and 111d are formed by NMOS transistors that the polarity thereof is the same as that of the first semiconductor switches 104a, 104b, 104c and 104d. \nHerein, the second semiconductor switch 111a of the memory block 106a, for instance, is connected at the source thereof to the higher order source line 110a while being connected at the drain thereof to the lower order source line 112a, and is also connected at the gate thereof to a second selected gate line 113a shared with another memory block 106b aligned along the row direction in the upper part. Thus, a predetermined gate voltage can be configured to be equally applied from the shared, single second selected gate line 113a to the two second semiconductor switches 111a and 111b that are mounted on the different memory blocks 106a and 106b disposed in the upper part.\nOn the other hand, similarly also in the two memory blocks 106c and 106d aligned along the row direction in the lower part, a single second selected gate line 113b is connected to the two second semiconductor switches 111c and 111d, and a predetermined gate voltage can be configured to be equally applied to the two second semiconductor switches 111c and 111d from the shared, second selected gate line 113b. \nFurther, each memory cell transistor 103 on the memory block 106a, for instance, is connected at one terminal thereof to the lower order bit line 105a while being connected at the other terminal thereof to the lower order source line 112a, and thus, the memory cell transistors 103 are disposed in parallel to each other between the lower order bit line 105a and the lower order source line 112a. The word lines 102a, 102b, 102c and 102d, shared by the memory block 106a and another memory block 106b aligned along the row direction, are connected to control gates of the memory cell transistors 103 of the memory block 106a. Due to this, a predetermined gate voltage can be configured to be equally applied, for instance, from the shared, single word line 102a to one of the memory cell transistors 103 of the memory block 106a in the upper part and one of the memory cell transistors 103 of another memory block 106b aligned with the memory block 106a along the row direction in the upper part.\nIncidentally, all the memory cell transistors 103 have the same structure that a channel region is disposed between one terminal and the other terminal, which are formed at a predetermined interval on a semiconductor substrate, and a charge accumulation layer, an interlayer insulation layer and a control gate are sequentially laminated through a tunnel insulation layer on the channel region of the semiconductor substrate. Such memory cell transistors 103 are of an N-channel type, and are capable of executing either data writing by injecting electric charges into the charge accumulation layer or data erasing by extracting the electric charges accumulated into the charge accumulation layer by means of voltage to be applied to the control gate and the region between the one terminal and the other terminal.\nThe non-volatile semiconductor memory device 100 thus structured can be configured to write data in a predetermined one of the memory cell transistors 103, read data from a predetermined one of the memory cell transistors 103 and erase data written in the memory cell transistors 103 by regulating voltages to be respectively applied to the higher order bit lines 101a and 101b, the higher order source lines 110a and 110b and the word lines 102a to 102h and by controlling on/off states of the first semiconductor switches 104a to 104d and the second semiconductor switches 111a to 111d. \nIn FIG. 6, the memory cell transistor 103 in the first row of the memory block 106a is set as a selected memory cell transistor 115 in which data is written, whereas all the remaining memory cell transistors 103 are set as non-selected memory cell transistors 116 in which data is not written.\nIt should be herein noted that for the sake of explanatory convenience, the memory block 106a on which the selected memory cell transistor 115 is disposed will be referred to as a selected block 117, whereas the memory blocks 106b, 106c and 106d on which only the non-selected memory cell transistors 116 are disposed will be referred to as non-selected blocks 118.\nActually in the non-volatile semiconductor memory device 100, for instance, when data is written in only the selected memory cell transistor 115 in the first row of the selected block 117, a high voltage of 12[V] is applied to a word line 120 (hereinafter referred to as a selected word line) that is the one connected to the selected memory cell transistor 115 among the plural word lines 102a to 102h, whereas a low voltage of 4[V] is applied to the word lines 121 (hereinafter referred to as non-selected word lines) that are the other remaining ones among the plural word lines 102a to 102h. \nFurther, at this time, in the non-volatile semiconductor memory device 100, a low voltage of 0[V] as a writing voltage can be applied to a higher order bit line 122 (herein referred to as a selected bit line) that is the one to which the selected memory cell transistor 115 is connected, whereas a high voltage of 8[V] as a writing prevention voltage can be applied to a higher order bit line 123 (herein referred to as a non-selected bit line) that is the one to which only the non-selected memory cell transistors 116 are connected. Moreover, in the non-volatile semiconductor memory device 100, a gate voltage of 10[V], which is higher than the voltage of the non-selected bit line 123, can be applied to the first semiconductor switches 104a and 104b from the first selected gate line 108a connected to the selected block 117, whereas a gate voltage of 0[V] can be applied to the second semiconductor switches 111a and 111b from the second selected gate line 113a. \nAccordingly, in the non-volatile semiconductor memory device 100, the first semiconductor switch 104b on the non-selected bit line 123 is switched on by means of the writing prevention voltage from the non-selected bit line 123 and the gate voltage from the first selected gate line 108a, and a writing prevention voltage of 8[V] can be applied to the non-selected memory cell transistor 116 on the non-selected bit line 123 that intersects with the selected word line 120. At this time, the second semiconductor switches 111a, 111b, 111c and 111d are switched off by applying a voltage of 0[V] thereto from the higher order source lines 110a and 110b and by applying a voltage of 0[V] thereto from the second selected gate lines 113a and 113b, and the lower order source lines 112a, 112b, 112c and 112d can be turned into a floating state.\nThus, in the non-selected memory cell transistor 116 at which the selected word line 120 and the non-selected bit line 123 intersect with each other, a voltage difference is reduced between the control gate and the semiconductor substrate, and as a result, electric charges cannot be injected into the charge accumulation layer without occurrence of a quantum tunneling effect.\nFurther, at this time, the first semiconductor switch 104a on the selected bit line 122 is switched on by means of the writing voltage from the selected bit line 122 and the gate voltage from the first selected gate line 108a, and a writing voltage of 0[V] can be applied to the selected memory cell transistor 115 on the selected bit line 122 that intersects with the selected word line 120. Accordingly, in the selected memory cell transistor 115 at which the selected word line 120 and the selected bit line 122 intersect with each other, a voltage difference is increased between the control gate and the semiconductor substrate by means of the writing gate voltage applied from the selected word line 120. As a result, a quantum tunneling effect occurs and electric charges can be injected into only the charge accumulation layer of the relevant selected memory cell transistor 115. Consequently, in the non-volatile semiconductor memory device 100, only the selected memory cell transistor 115 can be set to be in a data written state while electric charges are accumulated into the charge accumulation layer thereof."}
-{"text": "1. Field of the Invention\nThe present invention relates to a direct coupling type power amplifier for use with audio circuits and more particularly to such a direct coupling type power amplifier having a muting function.\n2. Description of the Prior Art\nIn a typical type of audio power amplifier circuit for amplifying audio signals, an inpact or popping sound unpleasant to the ears is sometimes generated from its load speaker immediately after the power source is turned on, due to a rapid change in bias voltages in a preamplifier such as an equalizer used in a first input amplifier stage of the power amplifier circuit. The impact sound is not only unpleasant to the ears, but also exposes the speaker to the danger of being damaged.\nThere has been provided a power amplifier circuit free from the generation of impact or popping sounds, in which a muting function is provided by breaking a signal transfer path through the actuation of an electromagnetic relay for a certain time interval after the turn-on of the power source. However, such a power amplifier circuit is expensive since the electromagnetic relay used therein is rather costly and moreover the switching sound of the relay is felt to be unpleasant."}
-{"text": "Services such as Google Maps (maps.google.com) allow users to conduct searches that are limited to particular geographic areas. For example, while viewing a street or satellite map of New York City on a web browser such as Google Chrome, a user may search for the word \u201cpizza\u201d. In response, the servers cause the user's computer to display icons indicating the positions of listings that match the query, such as restaurants that have \u201cpizza\u201d in their name.\nThe results of the query may be sent in a way that permits a great deal of interactivity. For example, the map is sent as tiles of images and descriptions of the results and their positions are sent as text (such as data formatted as JSON (Javascript Object Notation) or HTML). This allows the user's computer to display icons on the map and a text listing of the results next to the map. Using Javascript provided by the server, the client computer can determine whether the user has clicked an area of the map occupied by an interactive icon and perform a number of actions in response, such as displaying a balloon containing even more information about the listing.\nGoogle Maps also permits layers. For example, it is capable of showing icons for entries on Wikipedia.org that are associated with latitude/longitude coordinates matching the geographic area currently being viewed (e.g., http://en.wikipedia.org/wiki/Wikipedia:WikiProject_Geographica 1_coordinates). The icons of entries are sent as an image to be overlaid on the map, along with information relating to the entries. The information is sent in JSON with information identifying the entry's position on the map and an identification of the entry. The position information may be used by the user's computer to determine whether the user actuated the icon (such as by clicking it or hovering a mouse cursor over it) and the identification information may be used to obtain more information from a server if the icon is actuated. YouTube results may also be shown as a layer. These results are not filtered based on the query but, rather, show all entries in the area covered by the map.\nGoogle Maps also features a \u201cPlaces of Interest\u201d layer that allows users to display icons associated with all businesses falling within certain pre-selected categories sent by the server, such as \u201cItalian\u201d and \u201cPizza\u201d restaurants, banks, gas stations, and the like. Although the user may collectively display search results from the pre-selected categories and from searches based on custom search terms that are not pre-selected, the results are separate, i.e., the preselected category icons do not change based on the user-defined term and vice versa. For example, the user may check a box next to \u201cpizza\u201d and type in the term \u201cvegetarian\u201d, in which case Google Maps combines and displays the search results together; however, the results associated with \u201cpizza\u201d are not selected based on the \u201cvegetarian\u201d query nor are the results associated with the \u201cvegetarian\u201d query selected based on the selection of \u201cpizza\u201d from the list. Sites such as www.tavernoxoros.gr (showing many locations on a map in response to a user's selection of pre-selected categories) and www.smalltownbrooklyn.com (which may be capable of showing all businesses on a map) are similarly able to show all results associated with pre-selected categories."}
-{"text": "Activity in both humans and even animals depends largely on functioning, active muscles. As a consequence, sensing and monitoring of muscle activity is important for many reasons. For instance, in sports medicine, information about the fatigue level of a muscle is useful for preventing over-exertion of muscles, which might lead to muscle-based injuries. In addition, sensing muscle activity could be leveraged for precision training, allowing athletes to monitor progress on targeted muscle groups as needed. Finally, other ubiquitous applications such as muscle-controlled prosthetics could also be enabled by muscle activity data.\nThere are two main methods for detecting and monitoring muscle activity. The first is electromyography (EMG), which accurately measures the electrical potential of active muscle fibers. By amplifying and processing these electrical signals, information about muscle activation and activity can be inferred. The second method is mechanomyography (MMG). Mechanomyographic methods use vibration transducers to detect skeletal muscle activity. A brief discussion of the physiology of skeletal muscle vibrations that enable mechanomyography based systems is provided herein.\nFIG. 1 illustrates an example skeletal muscle 100 and in particular the quadriceps and the tendons 106a, 106b and muscle belly 102 (including muscle fibers 104) are shown. Skeletal muscles 100 generally consist of three parts: a muscle belly 102, composed of muscle fibers 104, and two tendon ends 106a, 106b. Activating skeletal muscles during physical activity causes the muscle fibers 104 to contract, resulting in mechanical vibrations due to: tremors of the human motor system, muscle fibers 104 sliding against each other and artifacts detected when a muscle belly's 102 circumference increases during muscle contractions and decreases during relaxation. These skeletal muscle vibrations are carried through to and detected at the skin level. The vibrations vary depending on the muscle fatigue status, muscle group size, and location in the body e.g., skeletal muscle proximity to bones of different mass and size.\nAs a method for inferring muscle activity, EMG while accurate, has the major setback that it is difficult, invasive and even painful to use. It requires perfect contact with the skin, necessitating the use of gels in the case of surface EMG, or the insertion of small needles into the muscle, in the case of fine-wire EMG. Mechanomyography methods on the other hand, work by externally detecting the mechanical vibrations occurring in the muscle belly of skeletal muscles when the muscle is activated using vibration transducers. As such they are less invasive and more convenient for muscle monitoring use in a dynamic environment. Current vibration transducer technologies include fabric stretch sensors, force sensitive resistors (FSR) and a combination of any of or all of available inertial sensors including accelerometers, gyroscopes and magnetometers."}
-{"text": "An HVAC system generally includes a closed loop refrigeration system with at least one evaporator, at least one condenser and at least one compressor. As the refrigerant travels through the evaporator, it absorbs heat from a heat transfer fluid to be cooled and changes from a liquid to a vapor phase. After exiting the evaporator, the refrigerant proceeds to a compressor, then a condenser, then an expansion valve, and back to the evaporator, repeating the refrigeration cycle. The fluid to be cooled (e.g. air) passes through the evaporator in a separate fluid channel and is cooled by the evaporation of the refrigerant. The cooled fluid can then be sent to a distribution system for cooling the spaces to be conditioned, or it can be used for other refrigeration purposes.\nOne type of air conditioner system is a split system where there is an indoor unit or heat exchanger, which is generally the evaporator, and an outdoor unit or heat exchanger, which is generally the condenser. Often, the outdoor unit is placed outdoors and is subject to outdoor ambient conditions, particularly temperature. When the outdoor ambient temperature falls, the amount of heat being removed from the refrigerant in the condenser increases. The increased heat removal in the condenser can result in a decrease in the refrigerant pressure at the suction line to the compressor, commonly referred to as head pressure. The decrease in head pressure results in a lowering of the temperature of the refrigerant at the evaporator. When the temperature of the refrigerant at the evaporator becomes too low, icing of the system can occur. Icing is a condition when the temperature at the exterior of the evaporator is sufficiently low to freeze water present in the atmosphere. The ice formed by the water frozen on the surface reduces the available heat transfer surface and eventually prevents the proper operation of the HVAC system by inhibiting heat transfer and/or damaging system components.\nSome attempts to address the problem of icing have utilized the control of system pressure. In one approach, a variable speed condenser fan or a plurality of condenser fans having independent controls are used to control airflow over the condenser coil. As the amount of air passing over the coil decreases, the amount of heat transfer taking place at the coil decreases. Therefore, the temperature of the refrigerant in the condenser and the pressure of the system increase to allow the indoor coil to cool the air without icing problems. The use of the variable speed condenser fan or a plurality of condenser fans having independent controls has the drawback that it is expensive and requires complicated wiring and controls.\nAn alternate approach for the problem of low system pressure or icing is a parallel set of condensers in the refrigerant cycle, as described in U.S. Pat. No. 3,631,686. The parallel set of refrigerant condensers allows for two modes of operation. One mode of operation allows refrigerant to flow from only one of the refrigerant condensers. During this mode of operation, the condenser that does not permit the flow of refrigerant fills with liquid refrigerant. Because of this flooding, there is a reduction in the effective surface area of the condenser. The reduced surface area thereby reduces the ability of the condenser to remove heat from the refrigerant. Therefore, the temperature of the refrigerant in the condenser and the head pressure of the system increase allowing the indoor coil to cool the air without icing. The use of parallel refrigerant condensers has the drawback that it requires an additional condenser coil and additional piping, thereby increasing the space and cost required for installation. Another drawback associated with refrigerant flooding of the condenser coil is the resultant decrease in system capacity. Refrigerant normally available in a properly operating system is trapped in the condenser coil and not available to the compressor, thereby decreasing system capacity.\nAn additional alternate approach for the problem of low system pressure is the use of a valve that controls the discharge or flow of liquid refrigerant from the condenser to a receiver vessel downstream of the condenser to maintain control of the amount of condensing surface exposed to the outside temperature, as described in U.S. Pat. No. 2,874,550. The discharge of refrigerant from the condenser is controlled by a pressure-response valve that mechanically opens to allow the flow of liquid refrigerant from the condenser to the receiver vessel reducing the level of liquid inside the condenser, thereby lowering the system pressure. Alternatively, the valve is closed to stop the flow until the level of refrigerant rises in the condenser in an amount that reduces the effective cooling surface of the condenser. The reduced surface area thereby reduces the ability of the condenser to remove heat from the refrigerant, thereby raising the pressure of the system. The use of a pressure-response valve and a vessel downstream of the condenser to maintain control of the amount of condensing surface has the drawback that it includes a specially designed valve and additional piping, thereby increasing the required space and cost. As discussed above, another one of the drawbacks with refrigerant flooding the condenser coil is decreased system capacity. Refrigerant normally available in a properly operating system is trapped in the condenser coil and not available to the compressor, thereby decreasing system capacity.\nAn additional alternate approach for the problem of low system pressure is the use of a refrigerant bypass around the condenser, as described in U.S. Pat. No. 3,060,699 and U.S. Reissued Pat. No. Re. 27,522. If the temperature and pressure of the refrigerant in the condenser are sufficiently high, a valve will close on a condenser bypass and the flow of refrigerant will be directed to the condenser. If the temperature and pressure of the condenser are not sufficiently high, the valve will open on a condenser bypass and at least some of the flow of refrigerant will be directed away from the condenser. The result of the bypass is an increase in pressure through the pipe leading to the evaporator downstream of the compressor. The use of a bypass has the drawback that it includes a specially designed valve and additional piping, thereby increasing the required space and cost.\nWhat is needed is a method and system for controlling the system refrigerant pressure without the drawbacks discussed above."}
-{"text": "This application is based upon and claims benefit of priority of Japanese Patent Application No. Hei-11-366454 filed on Dec. 24, 1999, the content of which is incorporated herein by reference.\n1. Field of the Invention\nThe present invention relates to a method of manufacturing a semiconductor device on an SOI (Silicon on Insulator) substrate, and more particularly to a method of forming an oxidized film on a silicon layer of the SOI substrate.\n2. Description of Related Art\nIn a conventional method of manufacturing capacitors or MOS transistors on a silicon substrate, an impurity is diffused into the silicon substrate in a high density to form a highly doped layer, and then an oxidized film is formed on the highly doped layer. Electrodes made of poly-silicon or the like are formed on the silicon dioxide film.\nIt is frequently required that such a silicon dioxide film has a high quality to be durable against a high voltage. However, such a high quality is difficult to be realized when the oxidized film is formed on the highly doped layer. This is because the highly doped layer tends to become a gettering site and the silicon dioxide film is contaminated with contaminants such as heavy metals, and the durability of the silicon dioxide film is adversely affected by the contaminants.\nIt is possible to form an additional gettering layer in the silicon substrate to improve the quality of the oxidized film. However, it is generally difficult to form such a gettering layer in the SOI substrate. JP-A-10-25621 proposes to form the gettering layer in the SOI substrate, but at least an additional costly process to form the same is required in a series of manufacturing processes of a semiconductor device.\nThe present invention has been made in view of the above-mentioned problem, and an object of the present invention is to provide a method of forming a quality oxidized film on the SOI substrate without using the process of forming the gettering layer.\nAn impurity such as boron or phosphorus is diffused into a silicon layer on an SCI substrate by ion-implantation and driven-in, thereby forming a diffused layer having a high density of the impurity on the SOI substrate. Then, an oxidized film is formed on the diffused layer by thermal oxidation. The oxidized film is heat-treated at a high temperature of 1100xc2x0 C. or higher, preferably 1150xc2x0 C. or higher. When boron is diffused to form the diffused layer, the high temperature heat treatment is performed for 20 minutes or longer. When phosphorus is diffused, the high temperature heat treatment is performed for a period of 20 minutes or longer but not exceeding 60 minutes.\nBy applying the high temperature heat treatment to the oxidized film, contamination in the oxidized film is driven-out, or annealed-out. Thus, the quality of the oxidized film such as durability against a high voltage is improved.\nTo form a capacitor using the oxidized film, electrodes made of a material such as poly-silicon are patterned on the oxidized film. Other semiconductor elements such as an LDMOS (Lateral Diffused MOS) may be formed on the same SOI substrate in addition to the capacitor. In this case, the high temperature heat treatment can be performed as a common process for improving the quality of the oxidized film and for forming other semiconductor elements.\nFurther, plural capacitors, e.g., a first capacitor and a second capacitor may be formed using the oxidized films having respectively different thickness. In this case, a preliminary oxidized film covering both areas for the first and second capacitors is first formed. Then, only the oxidized film covering the second capacitor area is removed by wet-etching, keeping the preliminary oxidized film for the first capacitor. Then, the silicon layer is again thermally oxidized to form a new oxidized film covering both areas for the first and second capacitors. Thus, a thin film is formed for the second capacitor and a thick film is formed for the first capacitor. The high temperature heat treatment is performed either immediately after both oxidized films are formed or after the capacitor electrodes are formed on the oxidized films. The quality of both oxidized films is improved by a single high temperature heat treatment.\nAccording to the present invention, the quality of the oxidized film or films formed on the SOI substrate is improved by the high temperature heat treatment without forming an additional gettering layer.\nOther objects and features of the present invention will become more readily apparent from a better understanding of the preferred embodiment described below with reference to the following drawings."}
-{"text": "This invention relates to a process for the melt-spinning of a filamentary structure from a synthetic polyamide polymer. More particularly, it is concerned with a continuous process for the formation of high strength polycaproamide filament, yarn or the like by melt-spinning a synthetic linear fiber-forming polycaproamide having excess number of carboxyl end groups over amino end groups.\nU.S. Pat. No. 2,174,527 discloses that dibasic dicarboxylic acids when employed in excess in a diacid-diamine polyamide, serve to terminate the polymer, thereby minimizing further increase in polymer molecular weight.\nU.S. Pat. No. 2,241,321 teaches polymerization of caprolactam in the presence of a diamine, followed by addition of sebacic acid and completion of the polymerization. the sebacic acid apparently acts as a chain-extending agent.\nMore recently, U.S. Pat. No. 3,386,967 discloses that a high strength polycaproamide yarn having an excess number of carboxyl end groups over amino end groups may be produced from a high molecular weight polymer prepared by polymerizing e-caprolactam and reacting the polymer with a dibasic carboxylic acid containing at least six carbon atoms. Because polycondensation of the polymer in the presence of a dibasic carboxylic acid is relatively slow, the acid is desirably included at the start of the polymerization reaction. Even so, equilibrium is approached in the polymerization reaction mixture at number average molecular weight not above about 20,000 under usual polymerization conditions, and special measures are required to carry the reaction further. One particularly useful method of accomplishing the required further reaction is to remove volatile by-products of the polymerization such as water by flowing an inert gas across the reaction mixture surface, desirably followed by application of vacuum.\nAlthough the process of U.S. Pat. No. 3,386,967 constitutes a major contribution to this art, we have found that in continuous operation of the process on a commercial basis, particularly at maximum production rates, serious problems have been encountered in melt spinning due to the frequent occurrence of \"nubs\" in the fiber. The term \"nubs\" is conventionally applied and is used herein to mean enlarged sections of filament no more than several filament diameters in length. Nubs may be formed by a foreign, nonorientable substance which interferes with normal fiber stretch in a short section, resulting in an enlargement. Foreign substances which are believed to have contributed to nuts in the present instance include carbonized polymer and gels formed in the polymer. Gels appear to be the chief cause, i.e., the nubs are probably created by non-orientable gel from cross-linked polymer. Thermal degradation of the polymer may be an important causative factor.\nThe reactions in thermal degradation of polyamides containing dicarboxylic acid additives are not entirely understood. It is likely that thermal degradation produces a decomposition product which serves to form cross-links between amide groups and adjacent polymer chains. The decomposition reaction proceeds slowly, finally building up a three-dimensional network of molecules which may be called polymer gel and which eventually reaches the stage where it forms an infusible coating on the walls of the equipment.\nA serious difficulty which arises from the formation of this polymer gel on the interior walls is that from time to time pieces break off and get into the flowing polymer stream where they produce damage to the spinning equipment.\nThe greatest difficulty, however, is caused by polymer gel which has progressed to the three-dimensional structural stage, but which has not yet reached the stage of being infusible. This kind of polymer gel is readily carried with the stream of flowing polymer. Being still molten or at least softened, it passes through the pump and even through the filter medium to show up either as discontinuities or as viscosity differences in the spun filament. When these filaments are later drawn, these defects may cause breaks in the filaments which either cause the whole thread to break or else form nubs which go through to be counted as quality defects in the final yarns."}
-{"text": "People routinely undergo surgery for treatment of a variety of medical conditions. In many cases, the patients are placed under general anesthesia for the surgical procedure. While under general anesthesia, the patient is unconscious and cannot respond or react to external stimuli or conditions. In some surgical procedures, there is a potential for these external conditions to cause trauma to the patient. For example, during certain orthopaedic surgical procedures, the surgeon must exert considerable pressure on various parts of the body, e.g., to the arms, legs, shoulder, and torso. In addition, the surgeon uses a variety of instruments during the procedure. The pressure on the body and use of surgical instruments create the potential for unintentional trauma to the patient, particularly to the head, face, ears, eyes, and neck area. The patient may receive pressure marks or sores from direct contact with surgical instruments and body positioning techniques. Any pressure applied to the ocular structures, even a short time, can cause damage or blindness to the eye. The body is typically covered at most by a blanket or cloth, which provides minimal protection.\nAnother challenge is the task of rolling the patient over from a supine position to a prone position on the operating table or from a cart onto the operating table. Depending on the procedure, the patient may be intubated, which creates a risk of neck injury during the roll-over process. After the procedure, the patient must again be rolled off the operating table onto a gurney. Still anesthetized, the risk of neck injury is again present if the head is not properly supported and manipulated during patient movement.\nIf an emergency develops while the patient is in the prone position, requiring the patient to be rolled to the supine position, valuable time can be lost trying to properly support the patient without injury to the neck, and without crimping the airway supply tubing and monitoring equipment communicating through the nose and mouth of the patient.\nOne approach found in U.S. Pat. No. 6,490,737 involves a molded helmet, which partially covers the patient's face and head. The helmet is particularly designed for procedures with the patient in the prone position and, in fact, the helmet can be rigidly mounted to the operating table, as shown in FIG. 8 of U.S. Pat. No. 6,490,737. However, in the prone position, the forehead and chin areas are exposed to continuous pressure by the weight of the patient's own head. If not relieved by regular movement of the face to allow blood flow, the pressure can cause localized ischemia to the chin and forehead area. In addition, the helmet leaves the eyes, cheeks, nose, and back of the head exposed and vulnerable to objects smaller than the openings in the helmet. The helmet is molded to a rigid form factor so the head is likely to shift in position relative to the helmet.\nIn the anesthetized state, the patient is unable to respond or react to any of these conditions. The surgeon, anesthesiologist, and surgical staff must be constantly aware of patient safety to avoid unnecessary trauma.\nThe operating room is typically maintained at a low temperature, in part for the comfort and alertness of the surgical team. The patient may experience body heat loss in the low temperature environment. The torso and upper and lower extremities are typically covered by a blanket for warmth. However, a significant amount of heat can be lost passively through the head and neck. A surgical cap can reduce some heat loss, but typically does not cover the face and certainly does not remove the potential for injury."}
-{"text": "The use of a CATV system to provide internet, voice over internet protocol (\u201cVOIP\u201d) telephone, television, security, and music services is well known in the art. In providing these services, a downstream bandwidth (i.e., radio frequency (\u201cRF\u201d) signals, digital signals, and/or optical signals) is passed from a supplier of the services to a user, and an upstream bandwidth (i.e., radio frequency (\u201cRF\u201d) signals, digital signals, and/or optical signals) is passed from the user to the supplier. For much of the distance between the supplier and the user, the downstream bandwidth and the upstream bandwidth make up a total bandwidth that is passed via a signal transmission line, such as a coaxial cable. The downstream bandwidth is, for example, signals that are relatively higher in frequency within the total bandwidth of the CATV system while the upstream bandwidth is, for example, signals that are relatively lower in frequency.\nTraditionally, the CATV system includes a head end facility, where the downstream bandwidth is initiated into a main CATV distribution system, which typically includes a plurality of trunk lines, each serving at least one local distribution network. In turn, the downstream bandwidth is passed to a relatively small number (e.g., approximately 100 to 500) of users associated with a particular local distribution network. Devices, such as high-pass filters, are positioned at various points within the CATV system to ensure the orderly flow of downstream bandwidth from the head end facility, through the trunk lines, through the local distribution networks, and ultimately to the users.\nAt various locations between the head end facility and the user, there are amplifiers and slope adjustment devices for the purpose of maintaining the quality of the downstream bandwidth. This statement introduces three terms (i.e., quality, amplifiers, and slope adjustment devices) that are important to the remaining discussion. These will be discussed broadly below.\nThe quality of the downstream bandwidth is often a measure of: (i) a signal level of a particular channel within the downstream bandwidth, the signal level referred to merely as \u201clevel;\u201d and (ii) a general consistency of levels across all of the channels in the downstream bandwidth, the general consistency referred to as \u201cslope.\u201d These objective measurements are often used by technicians to evaluate CATV system performance during operation and to troubleshoot customer complaints.\nThe level of each channel should fall within a specific range that has been determined to provide satisfactory video, sound and information transfer rates for users. The specific requirements for each channel are not of importance to the present discussion, but it is helpful to understand that are specific targets for the level of each channel. Note that this is a simplistic definition to explain \u201clevel,\u201d and note that this definition does not include other variances such as between analog and digital.\nSlope is measurement used to assess the amount of loss experienced due in large part to cable length. While all channels experience some loss, channels transmitted using higher frequencies within the downstream bandwidth experience more loss than those transmitted using lower frequencies. Accordingly, when the levels for all of the channels within the downstream bandwidth are graphed such that they are arranged in order according to the frequency of the channel, there may be a significant visual downward slope in the graph from the lowest frequency channel to highest frequency channel. This downward slope becomes more prominent as the length of signal cable increases. Note that this is a simplistic definition to explain the consistency of levels across all of the channels and the \u201cslope\u201d that is created by losses occurring in the signal cables. Also note that this definition does not include other variances such as between analog and digital.\nThe presence of slope is not removed through the use of typical drop-style amplifiers. The drop-style amplifiers merely amplify the entire downstream bandwidth. In other words, these drop-style amplifiers raise the level of each channel equally. In turn, if there is a large amount of slope present, such as when a user's premise includes long distances of signal cable, the drop-style amplifier may cause some channels to exceed their level specification while other channels may remain below their specification.\nIt is known to add a fixed or manually adjustable slope compensator/low frequency attenuator when there is a long run of signal cable. However, these devices require expensive testing equipment to determine whether and/or how much slope compensation should be supplied to a particular premise. Further, due to the cost of installation and a general misunderstanding regarding how to install such devices, there are relatively few in existence, compared to the number of such devices needed.\nFor at least the forgoing reasons, a need is apparent for a device, which can increase the overall quality of the downstream bandwidth that includes providing an appropriate amount of amplification and an appropriate amount of slope compensation increasing the signal strength and increasing the signal-to-noise ratio."}
-{"text": "1. Field of the Invention\nThe present invention generally relates to a program register using a nonvolatile memory device and a programmable logic circuit using the same, and more specifically, to a technology for storing data or performing an operation on the data without additional memory devices, thereby reducing the area of the circuit.\n2. Description of the Prior Art\nGenerally, a ferroelectric random access memory (hereinafter, referred to as \u2018FRAM\u2019) has attracted considerable attention as next generation memory device because it has a data processing speed as fast as a Dynamic Random Access Memory DRAM and conserves data even after the power is turned off.\nThe FRAM having structures similar to the DRAM includes the capacitors made of a ferroelectric substance, so that it utilizes the characteristic of a high residual polarization of the ferroelectric substance in which data is not deleted even after an electric field is eliminated.\nThe technical contents on the above FRAM are disclosed in the Korean Patent Application No. 1999-49972 by the same inventor of the present invention. Therefore, the basic structure and the operation on the FRAM are not described herein.\nA conventional programmable logic operation circuit for changing logic levels of input signals stores address information in storage means. However, since a SRAM (Static Random Access Memory) is used as the conventional programmable logic operation circuit, various information stored in latches is leaked in a power-off mode. Even when power is supplied to the system again, various data for operations of circuits are required to be reset."}
-{"text": "In a method for detecting analyte which is well-known in this field, latex or erythrocyte is sensitized with antigen or antibody beforehand, and, then, agglutination reaction with an analyte which exists in biological fluid such as protein or a certain specific organic material is utilized for detecting the analyte. An example of latex which has widely been employed for the above-mentioned method is polystyrene latex, which is being studied also as to what influence is exerted on analyte-detecting rate or the sensitivity of detection (e.g., non-specific adsorption of materials other than analyte, or the like) by factors such as particle size and functional group (carboxyl group and the like) on the surface of particle (see Journal of the Japan Society for Clinical Laboratory Automation, 12, 2, 121-124, 1987).\nAgglutination of latex particles via analyte has a major influence on detecting rate. It takes, however, a considerable time for the agglutination, since the speed of said agglutination is determined by diffusion. Hence, it has been tried, for instance, to make latex particles carry magnetic particles and to raise the rate of agglutination or precipitation of latex particles by magnetic force (see, e.g., J. Applied Polymer Science, Vol. 50, 765-776, 1993).\nPolystyrene latex particles and the like which are mentioned in the former of the above-mentioned documents are liable to agglutinate non-specifically, and are sometimes slightly inferior in dispersion stability. Furthermore, the surface of such particles is apt to non-specifically adsorb impure protein other than analyte. Thus, a part of the inventors of this invention has provided, as a quantum dot to detect analyte in biological sample, composite particles which are composed, for instance, of semiconductor-ultra fine particles-encapsulating fine particles whose surface is covered with chain of polyethylene glycol (hereinafter referred to also as PEG) or poly(ethylene oxide) (or wherein PEG having high mobility in an aqueous solution takes plural brush-like structures) (see pamphlet of WO 02/056020, pages 16-17). Although this type of quantum dot prevents non-specific adsorption of impure protein, it has generally a particle size up to hundreds of nanometers, and has high-mobility PEG on its surface. Thus, quantum dot makes a very stable dispersion in an aqueous solution, and has therefore not necessarily been considered to be suitable for use in a system to rapidly detect analyte with use of agglutination reaction."}
-{"text": "The present invention relates to a technique for adding an inherent item, which is to be referred to or to be updated by a particular organization, to information (such as a table) shared among a plurality of organizations in a Relational Data Base Management System (RDBMS).\nThe RDBMS is a system to manage data shared among a plurality of organizations. For example, in an information system of a firm or a company, corporate information of firms as its customer is stored in the RDBMS to share the corporate information among organizations (e.g., divisions) of the firm. The corporate information is stored in the RDBMS in the form of a table including a plurality of items (columns). A business application (AP; program) accesses the table to refer to or to update the corporate information. The items are, for example, names and locations of the firms.\nIn such data processing system, the respective organizations desire to use mutually different items depending on cases. FIG. 2 shows an example of this situation in a customer corporate information management system. In the customer corporate information management system 100, in-house organizations such as a sales division and a supply division share customer corporate information 130. However, the respective divisions desire to use different items. A customer management application (sales division customer management AP 110) of the sales division refers to information 111 to use items which store a salesperson and a contact address, e.g., a contact phone number (salesperson contact phone number) of the salesperson of an associated firm. On the other hand, an order application (supply division order AP 120) of the supply division refers to information 121 to use a state of approval and an approval number of International Organization for Standardization (ISO) 9001, i.e., the international standard for quality certification, to thereby determine whether or not a firm under consideration is appropriate as a supplier.\nAssume that items shared among the organizations are referred to as shared items and items used only by a particular organization are referred to as inherent items. In the example of FIG. 2, the firm names and the firm locations are shared items. The salesperson and the salesperson contact phone number are inherent items used only by the sales division.\nIt is not assumed that the one-to-one correspondence exists between the business applications (sales division customer management AP 110 and supply division order AP 120 in this example) and the organizations. Depending on cases, a business application is used by only one organization or by a plurality of organizations. In the former situation, the business application uses inherent items of a particular organization. In the latter case, the business application uses inherent items of organizations to which the users of the application belong. That is, even one business application uses different inherent items depending on users.\nIn the conventional RDBMS, the inherent items are managed in three methods as below.\nAccording to a first method, for each organization, a table is prepared to store inherent items (which will be described in conjunction with FIG. 25). In a second method, for each organization, a table is prepared to store both of shared items and inherent items (which will be described in conjunction with FIG. 26). In a third method, a table is prepared to store both of shared items and inherent items for the respective organizations (which will be described in conjunction with FIG. 27).\nThese three methods have been introduced as three mapping patterns provided by Hibernate which is an object/relational mapping tool in wide use, in Section 3.6 \u201cClass Inheritance Mapping\u201d of \u201cHibernate In Action\u201d. Reference is to be made to pages 121 to 131 of \u201cHibernate In Action\u201d written by Christian Bauer et al, translated by Akira Kurahashi et al, and published from Softbank Creative on Jan. 15, 2006.\nMethods 1 to 3 respectively correspond to \u201c3.6.3 Table per subclass\u201d, \u201c3.6.1 Table per concrete class\u201d, and \u201c3.6.2 Table per class hierarchy\u201d of the article described above.\nIn general, according to the class structure (logical data structure) of objects to be mapped, shared items are designated as attributes of a main class and inherent items of each organization are designated as attributes of a subclass for each organization. FIG. 24 shows an example of the class structure for the data items of FIG. 2. The main class having a shared item as an attribute is associated with an in-house shared corporate attribute 2401, and the subclass having an inherent item as an attribute is associated with a sales division inherent corporate attribute 2402 and a supply division inherent corporate attribute 2403. The sales division inherent corporate attribute 2402 has, as attributes, inherent items to be used by the sales division. The supply division inherent corporate attribute 2403 has, as attributes, inherent items to be used by the supply division.\nDescription will be given of FIGS. 25 to 21 based on the above assumption. In FIG. 25, shared items are stored in an in-house shared table 2501, inherent items to be used by the sales division are stored in the sales division inherent table 2502, and inherent items to be used by the supply division are stored in the supply division inherent table 2503.\nIn FIG. 26, shared items 2602 and inherent items to be used by the sales division are stored in a sales division inherent table 2601, and shared items 2612 and inherent items to be used by the supply division are stored in the supply division inherent table 2611.\nIn FIG. 27, a shared item 2711, a sales division inherent item 2712, and a supply division inherent item 2713 are stored in one in-house shared table 2701."}
-{"text": "Polyolefin polymers are very well known in the art. Methods of manufacturing and/or processing polyolefins are likewise well known in the art. These methods of processing include but are not limited to blow molding, injection molding, pultrusion and pulforming, and extrusion. The problems associated with each of these processes are known, and seeking solutions to these problems, and/or improvements in processing conditions and equipment remains of interest to the art and industry.\nFor example, it is known that olefinic homopolymers such as high-density polyethylenes (HDPE) and polypropylenes generally have poor processability and poor impact properties. Because of the wide commercial uses of these polymers in various applications such as automotives, households, and. packaging, it is always desirable to discover new methods, additives, and/or equipment which enable these polymers to be better and more economically processed, and to improve the mechanical and other properties of these polymers and/or products made from the polymers."}
-{"text": "The market for portable devices, for example, mobile phones, smart watches, tablets, etc., is expanding with many more features and applications. As the number of applications on these devices increases, there also is an increasing demand to run multiple applications concurrently. More features and applications call for microprocessors to have high performance, but with low power consumption. Multithreading can contribute to high performance in this new realm of application. Keeping the power consumption for the microprocessor and related cores and integrated circuit chips near a minimum, given a set of performance requirements, is desirable, especially in portable device products.\nMultithreading is the ability to pursue two or more threads of control in parallel within a microprocessor pipeline. Multithreading is motivated by low utilization of the hardware resource in a microprocessor. In comparison, multi-core is fairly wasteful of the hardware resource. Multithreading can, in general, provide the same performance as multicore without duplicating of resources.\nMultithreading can be used in an effort to increase the utilization of microprocessor hardware and improve system performance. Multithreading is a process by which two or more independent programs, each called a \u201cthread,\u201d interleave execution in the same processor, which is not a simple problem. Each program or thread has its own register file, and context switching to another program or thread requires saving and restoring of data from a register file to a memory. This process can consume much time and power. These and other problems confront attempts in the art to provide efficient multithreading processors and methods."}
-{"text": "1. Field of the Invention\nThe invention relates to chemical vapor deposition (CVD) systems and, more particularly, to temperature control apparatus, and a concomitant method, for controlling the temperature of the walls of a CVD reaction chamber.\n2. Description of the Prior Art\nIn a conventional chemical vapor deposition (CVD) system, a reaction chamber envelopes a volume wherein the deposition process occurs. The chamber typically is fabricated of transparent quartz. As such, the reaction chamber is said to have quartz windows (accessible surfaces of the chamber) through which the temperature of the apparatus within the chamber can be optically monitored using optical pyrometers. In general, since the windows are integral with the walls of the chamber, the temperature of these windows is indicative of the chamber wall temperature. The temperature of the windows and the walls of the chamber vary depending upon the temperature of the workpiece within the chamber, the process presently being accomplished, the gases involved in the current process, and the run time of the process.\nMore specifically, it is desirable that during a deposition process, the quartz wall temperature be controlled within a narrow temperature range to minimize deposits on the chamber wall. Furthermore, it is desirable that, during cleaning steps, the quartz wall temperature be maintained at a different temperature to maximize the etch rate of any film that may have been deposited on the chamber walls. Such temperature maintenance minimizes particulate contamination problems.\nTypically, the temperature of the walls of the chamber is maintained at a nominal value by a flow of air proximate the outer surface of the chamber. This air flow is fixed for each process performed within the chamber. Thus, the temperature of the chamber walls varies with conditions within the chamber, e.g., the wall temperature differs depending on whether the process in the chamber is depositing material or etching material from a workpiece.\nThrough empirical study, it has been shown that the processes performed within the chamber can be optimized when the chamber wall temperature is maintained at an optimum temperature for a particular process. However, prior art CVD systems contemplate maintaining the air flow past the chamber at a fixed flow rate. As such, these systems exhibit significant chamber wall temperature variations which preclude a high degree of temperature optimization.\nTherefore, a need exists in the art for closed loop apparatus, and a concomitant method, for controlling the air flow proximate a reaction chamber such that the temperature of the chamber walls can be readily optimized for each of the processes accomplished within the chamber."}
-{"text": "Electromagnetic sources used in geophysical surveying generally produce unfocused regions of subsurface excitation. Unfocused excitations result both from the diffusive nature of low frequency electromagnetic wave propagation in the earth and from the use of relatively simple source configurations in the surveys. Unfocused excitations are a severe disadvantage when a target-oriented survey is needed, such as when it is desirable to probe the electrical properties of a prospective formation in a location where a complex geologic background may be present.\nTo eliminate the disadvantages of unfocused surveys, a method of source focusing is disclosed in U.S. patent application Ser. No. 09/656,191, filed Sep. 6, 2000 (hereinafter referred to as \u201cSrnka\u201d). The method, referred to as the remote reservoir resistivity mapping method (\u201cR3M\u201d) uses a spatially extended grounded electromagnetic source array that is positioned on the earth's surface, or near the seafloor, to focus electromagnetic energy on the subsurface target. The source array consists of two concentric electrode rings whose radii are optimally determined from theoretical modeling to focus the subsurface excitation at the approximate depth of the subsurface target. A limitation of the Srnka method, however, is that for typical reservoir depths, which range from 500 meters to greater than 5000 meters, the total length of electrode wire that is required in the concentric electrode rings ranges from about 23.5 km to greater than 235 km.\nThese extreme lengths of electrode wire raise numerous operational challenges. For example, difficulties in handling and deploying the electrode wires, and problems in ensuring accurate positioning of the wires in the desired electrode ring locations may be involved, and there is the need to avoid obstacles in the area of the desired ring locations, and the requirement that the electrodes be uniformly grounded along their lengths. Other problems influence the economics of the survey, including the large cost of the electrode wire, some of which is likely to be damaged or lost in each survey, and the large labor and logistics costs associated with mobilization, demobilization, deployment, and recovery of the wire.\nFor these reasons, Srnka also disclosed an alternative method that reduces the required length of electrode wire by substituting a set of equally spaced radial electric bipoles, each having a length equal to the difference in the lengths of the radii of the two concentric rings. FIG. 1 depicts both the R3M concentric ring method and this radial bipole alternative. Concentric electrode rings 4 and 5 lie on surface 1 of earth 2 and have radii a and b, respectively, and a center 7 generally above reservoir 3. In the radial bipole alternative method these concentric rings are replaced by radial electrode bipole sources 6. Each bipole source 6 has a length L=b\u2212a. Srnka discloses that a minimum of six radial bipole sources 6 is preferable. A signal equivalent to the signal that a receiver array would receive from the concentric ring source array is produced mathematically by summing the signals received by the receiver array from each of the bipole sources 6 (For convenience, FIG. 1 does not depict the location of the receiver array). This alternative method, which allows the simulation of a concentric ring source array, is referred to as the virtual concentric ring source method, or as the radial bipole method, and substantially reduces the required length of electrode wire as compared to the R3M physical concentric ring source method.\nThe virtual concentric ring source method nevertheless uses relatively long lengths of electrode wire and inherently retains many of the operational challenges of the R3M physical concentric ring source method. In addition, both methods have other inherent disadvantages, in particular for certain aspects of the subsurface resistivity-imaging problem. Specifically, both methods produce maximum responses near a pre-specified subsurface excitation focus depth (in other words, at the approximate depth of the target), but provide little differential sensitivity to features within an appreciable vertical range of that depth. More specifically, the ability to image separately features above or below the target a distance of less than five percent of the target's depth below the surface is limited. This limitation greatly hinders the imaging of stacked reservoirs. In addition, amplifier saturation effects can influence data quality and hinder analysis of such data. A third issue is that both methods produce primarily vertical subsurface currents in the region of the target. These currents are not adequate for imaging reservoir electrical macro-anisotropy (vertical electrical isotropy), such as exists in many reservoirs due to the presence of shale interbeds. Imaging this anisotropy is essential for estimating reservoir net and gross volumes, and for determining hydrocarbon pore volumes.\nAccordingly, a method is desired which retains the advantageous electromagnetic responses of the R3M physical concentric ring source and radial bipole methods, but which reduces or eliminates the disadvantages of those methods. The present invention satisfies that desire."}
-{"text": "The present invention relates to introduction of naked DNA and RNA sequences into a vertebrate to achieve controlled expression of a polypeptide. It is useful in gene therapy, vaccination, and any therapeutic situation in which a polypeptide should be administered to cells in vivo.\nCurrent research in gene therapy has focused on xe2x80x9cpermanentxe2x80x9d cures, in which DNA is integrated into the genome of the patient. Viral vectors are presently the most frequently used means for transforming the patient\"\"s cells and introducing DNA into the genome. In an indirect method, viral vectors, carrying new genetic information, are used to infect target cells removed from the body, and these cells are then re-implanted. Direct in vivo gene transfer into postnatal animals has been reported for formulations of DNA encapsulated in liposomes and DNA entrapped in proteoliposomes containing viral envelope receptor proteins (Nicolau et al., Proc. Natl. Acad Sci USA 80:1068-1072 (1983); Kaneda et al., Science 243:375-378 (1989); Mannino et al., Biotechniques 6:682-690 (1988). Positive results have also been described with calcium phosphate co-precipitated DNA (Benvenisty and Reshef Proc. Natl. Acad Sci USA 83:9551-9555 (1986)).\nThe clinical application of gene therapy, as well as the utilization of recombinant retrovirus vectors, has been delayed because of safety considerations. Integration of exogenous DNA into the genome of a cell can cause DNA damage and possible genetic changes in the recipient cell that could predispose to malignancy. A method which avoids these potential problems would be of significant benefit in making gene therapy safe and effective.\nVaccination with immunogenic proteins has eliminated or reduced the incidence of many diseases; however there are major difficulties in using proteins associated with other pathogens and disease states as immunogens. Many protein antigens are not intrinsically immunogenic. More often, they are not effective as vaccines because of the manner in which the immune system operates.\nThe immune system of vertebrates consists of several interacting components. The best characterized and most important parts are the humoral and cellular (cytolytic) branches. Humoral immunity involves antibodies, proteins which are secreted into the body fluids and which directly recognize an antigen. The cellular system, in contrast, relies on special cells which recognize and kill other cells which are producing foreign antigens. This basic functional division reflects two different strategies of immune defense. Humoral immunity is mainly directed at antigens which are exogenous to the animal whereas the cellular system responds to antigens which are actively synthesized within the animal.\nAntibody molecules, the effectors of humoral immunity, are secreted by special B lymphoid cells, B cells, in response to antigen. Antibodies can bind to and inactivate antigen directly (neutralizing antibodies) or activate other cells of the immune system to destroy the antigen.\nCellular immune recognition is mediated by a special class of lymphoid cells, the cytotoxic T cells. These cells do not recognize whole antigens but instead they respond to degraded peptide fragments thereof which appear on the surface of the target cell bound to proteins called class I major histocompatibility complex (MHC) molecules. Essentially all nucleated cells have class I molecules. It is believed that proteins produced within the cell are continually degraded to peptides as part of normal cellular metabolism. These fragments are bound to the MHC molecules and are transported to the cell surface. Thus the cellular immune system is constantly monitoring the spectra of proteins produced in all cells in the body and is poised to eliminate any cells producing foreign antigens.\nVaccination is the process of preparing an animal to respond to an antigen. Vaccination is more complex than immune recognition and involves not only B cells and cytotoxic T cells but other types of lymphoid cells as well. During vaccination, cells which recognize the antigen (B cells or cytotoxic T cells) are clonally expanded. In addition, the population of ancillary cells (helper T cells) specific for the antigen also increase. Vaccination also involves specialized antigen presenting cells which can process the antigen and display it in a form which can stimulate one of the two pathways.\nVaccination has changed little since the time of Louis Pasteur. A foreign antigen is introduced into an animal where it activates specific B cells by binding to surface immunoglobulins. It is also taken up by antigen processing cells, wherein it is degraded, and appears in fragments on the surface of these cells bound to Class II MHC molecules. Peptides bound to class II molecules are capable of stimulating the helper class of T cells. Both helper T cells and activated B cells are required to produce active humoral immunization. Cellular immunity is thought to be stimulated by a similar but poorly understood mechanism.\nThus two different and distinct pathways of antigen processing produce exogenous antigens bound to class II MHC molecules where they can stimulate T helper cells, as well as endogenous proteins degraded and bound to class I MHC molecules and recognized by the cytotoxic class of T cells.\nThere is little or no difference in the distribution of MHC molecules. Essentially all nucleated cells express class I molecules whereas class II MHC proteins are restricted to some few types of lymphoid cells.\nNormal vaccination schemes will always produce a humoral immune response. They may also provide cytotoxic immunity. The humoral system protects a vaccinated individual from subsequent challenge from a pathogen and can prevent the spread of an intracellular infection if the pathogen goes through an extracellular phase during its life cycle; however, it can do relatively little to eliminate intracellular pathogens. Cytotoxic immunity complements the humoral system by eliminating the infected cells. Thus effective vaccination should activate both types of immunity.\nA cytotoxic T cell response is necessary to remove intracellular pathogens such as viruses as well as malignant cells. It has proven difficult to present an exogenously administered antigen in adequate concentrations in conjunction with Class I molecules to assure an adequate response. This has severely hindered the development of vaccines against tumor-specific antigens (e.g., on breast or colon cancer cells), and against weakly immunogenic viral proteins (e.g., HIV, Herpes, non-A, non-B hepatitis, CMV and EBV).\nIt would be desirable to provide a cellular immune response alone in immunizing against agents such as viruses for which antibodies have been shown to enhance infectivity. It would also be useful to provide such a response against both chronic and latent viral infections and against malignant cells.\nThe use of synthetic peptide vaccines does not solve these problems because either the peptides do not readily associate with histocompatibility molecules, have a short serum half-life, are rapidly proteolyzed, or do not specifically localize to antigen-presenting monocytes and macrophages. At best, all exogenously administered antigens must compete with the universe of self-proteins for binding to antigen-presenting macrophages.\nMajor efforts have been mounted to elicit immune responses to poorly immunogenic viral proteins from the herpes viruses, non-A, non-B hepatitis, HIV, and the like. These pathogens are difficult and hazardous to propagate in vitro. As mentioned above, synthetic peptide vaccines corresponding to viral-encoded proteins have been made, but have severe pitfalls. Attempts have also been made to use vaccinia virus vectors to express proteins from other viruses. However, the results have been disappointing, since (a) recombinant vaccinia viruses may be rapidly eliminated from the circulation in already immune individuals, and (b) the administration of complex viral antigens, may induce a phenomenon known as xe2x80x9cantigenic competition,xe2x80x9d in which weakly immunogenic portions of the virus fail to elicit an immune response because they are out-competed by other more potent regions of the administered antigen.\nAnother major problem with protein or peptide vaccines is anaphylactic reaction which can occur when injections of antigen are repeated in efforts to produce a potent immune response. In this phenomenon, IgE antibodies formed in response to the antigen cause severe and sometimes fatal allergic reactions.\nAccordingly, there is a need for a method for invoking a safe and effective immune response to this type of protein or polypeptide. Moreover, there is a great need for a method that will associate these antigens with Class I histocompatibility antigens on the cell surface to elicit a cytotoxic T cell response, avoid anaphylaxis and proteolysis of the material in the serum, and facilitate localization of the material to monocytes and macrophages.\nA large number of disease states can benefit from the administration of therapeutic peptides. Such peptides include lymphokines, such as interleukin-2, tumor necrosis factor, and the interferons; growth factors, such as nerve growth factor, epidermal growth factor, and human growth hormone; tissue plasminogen activator; factor VIII:C; granulocyte-macrophage colony-stimulating factor; erythropoietin; insulin; calcitonin; thymidine kinase; and the like. Moreover, selective delivery of toxic peptides (such as ricin, diphtheria toxin, or cobra venom factor) to diseased or neoplastic cells can have major therapeutic benefits. Current peptide delivery systems suffer from significant problems, including the inability to effectively incorporate functional cell surface receptors onto cell membranes, and the necessity of systemically administering large quantities of the peptide (with resultant undesirable systemic side effects) in order to deliver a therapeutic amount of the peptide into or onto the target cell.\nThese above-described problems associated with gene therapy, immunization, and delivery of therapeutic peptides to cells are addressed by the present invention."}
-{"text": "1. Field of the Invention\nThis invention relates to an optical card handling system used with a general-purpose operating system and a method of accessing an optical card.\n2. Description of the Related Art\nInformation processing techniques have recently been making remarkable progress, and means for recording an increasingly large volume of information have been required. In this connection, what is now attracting attention is an optical information recording/reproducing apparatus.\nThis type of apparatus uses a recording medium which information can be optically written into and read from. As one of optical information recording/reproducing apparatuses of this type, an optical card apparatus using an optical card as a recording medium has been put to practical use. The optical card apparatus is capable of recording information on an optical card and reproducing the recorded information.\nThe optical card contains an information recording layer which has the property of changing irreversibly in the presence of heat. By projecting the laser light gathered by a lens onto the information recording layer, pits (holes) can be made in the layer. In this way, the presence and absence of pits (holes) made in the information recording layer represent a data string. Thus, by controlling the formation of pits according to the contents of data to be recorded, the desired data can be written. The reproduction of data is effected by optically sensing the presence and absence of pits in the information recording layer and converting into data.\nThe optical card has a recording capacity thousands to tens of thousands of times that of a conventional magnetic card. Once data is recorded in an area, the data cannot be rewritten. However, the data capacity that can be recorded on a single optical card is as large as 2 Mbyte to 4 Mbyte. Therefore, the optical card is considered to applied to a variety of uses including bankbooks, pocket maps, and prepaid cards used in shopping. Furthermore, by making use of the advantage that data cannot be rewritten, the optical card is considered to be applied to the case where data must not be falsified.\nAn example of the optical card is shown in FIG. 12. The optical card 1 shown in the figure is of the credit card size and has a data recording section 2 on its one side. The data recording section 2 is divided into a lot of tracks 2a as the magnetic recording medium is. Each track 2a is straight. Data is recorded sequentially track by track in the direction in which the track extends. The specific areas at both ends of each track 2a, that is, the specific regions at the begin and terminate ends of each track 2a, are used as ID sections 3a, 3b in which information items such as track addresses are recorded.\nAs described above, the optical card 1 is a medium on which data is recorded in the form of the presence and absence of pits, more specifically, on which data is recorded by forming pits in each track. Once pits are formed, they cannot be filled in, so that new data cannot be written over the track in which data is already written. This type of optical card is called a write-once, read-many type.\nHowever, data can be written in an unused track. Because each track has the ID sections 3a, 3b, a file can be updated as follows. When the data is updated, the whole file including the data items to be updated is recorded on an unused track, and the old data file is discarded. As with the magnetic disk, some tracks are used for file management. Once the file is updated, new file management information is written on unused ones among the file management tracks, and management is effected so that the updated file can be distinguished from the discarded file.\nTo record and reproduce data onto and from the optical card 1, a read/write beam of laser light must be swept track by track. The light beam may be swept along the track. In addition, the same result may be produced by placing the light beam at a constant position and causing the optical card to move repeatedly along the track.\nIn general, the optical card is moved along the track and the light source of the optical beam is moved across the track so that data can be read from and written into the desired track.\nA detailed structure of the data recording section 2 will be described. As shown in FIG. 13, a plurality of track guides 2a1 are provided at regular intervals in parallel with the track so that each track 2a2 may be formed between two adjacent track guides. Each track guide 2a1 can be sensed optically. A sensor for sensing the track guides is provided. Each track guide is sensed by the sensor. Because each track guide 2a1 and each track 2a2 are arranged along the same lane, the position of the light beam can be controlled so as not to depart from the track lane by sensing the track guide corresponding to the desired track with the sensor and controlling the position of the data read/write light beam so that the beam may move along the track guide. Of course, the data read/write light beam is positioned in the center of the track lane.\nIn this way, data pits 2a3 representing the recording information can be formed and read along the track 2a2. The track 2a2 of the data recording section 2 (not shown) is further divided into sectors. Data is read and written in sectors.\nThe data recording section 2 which data is read from and written into has a data area 2b1, a directory storage area 2b2, a change area 2b3, and a change information area 2b4. Namely, some of the tracks in the data recording section 2 are allocated to the data area 2b1, directory storage area 2b2, change area 2b3, and change information area 2b4, which are used as dedicated areas.\nThe change area 2b3 is an area used for changing the data when a recording error has occurred due to a defective track on the optical card 1, or for logically rewriting the data. The change information area 2b4 is an area in which the fact that the data is written in the change area 2b3 for replacement is recorded.\nThe data area 2b1 is an area in which the contents of the data forming a file is recorded. The directory storage area 2b2 is an area in which file management information is stored. In the directory storage area 2b2, the necessary information for file management, including a file name, the locations in which the file is recorded, and the length of the file, is recorded. Therefore, by searching the directory storage area 2b2 for the data storage location, the contents of the data can be retrieved.\nWith the optical card, the reading and writing of information is effected by an optical head composed of an optical system as shown in FIG. 15.\nIn the optical head shown in FIG. 15, the laser light emitted from a light-emitting element 4 is shaped by a collimate lens into parallel light and then passes through an object lens 6. The object lens 6 focuses the light onto the optical card 1. The reflected light from the optical card 1 passes through a mirror 7, an image forming lens 8, and a beam splitter 9, and is supplied to a detector 10. The detector 10 senses the reflected light and converts it into an electric light.\nThe detector 10 produces an electric signal according to the presence and absence of pits in the track. Thus, moving the optical card 1 along the track lane enables the data pits 2a3 shown in FIG. 13 to be sensed by the detector 10 in the form of the strength and weakness of the electric signal. Furthermore, in the optical head, the light obtained from the beam splitter 9 is supplied to a focus sensor 11 to sense the deviation of focus. By actuating an object lens driving section 12 according to the sense result of the focus deviation, the position of the object lens 6 is adjusted and thereby the focus is adjusted. To reciprocate the optical card 1 in its longitudinal direction (along the track lane), an optical card driving section 13 is provided. Explanation of a tracking error signal sensing system will be omitted.\nWhen the optical card 1 is used in a computer system, attention should be given to the following points. In the computer system, to use the optical card 1, an optical card recording/reproducing apparatus is connected via an interface. The optical card recording/reproducing apparatus is controlled in various ways under the control of the operating system (OS), the basic program for the computer system, and reads and writes data from and onto the optical card 1.\nMany personal computer systems and workstations use UNIX (a trademark of AT&T), MS-DOS (a trademark of Microsoft Corporation), Windows (a trademark of Microsoft Corporation), Macintosh-OS (a trademark of Apple Computer, Inc.), and DOS/V (a trademark of IBM corporation) as an operating system (OS). Under the present conditions, these OSes occupy the position of standard OSes in the fields of personal computer systems and workstations.\nWith these standard OSes, to manage files for an external storage apparatus, various pieces of information on files are stored in the directory, and on the basis of the information in the directory, the reading and writing of files is effected. These standard OSes differ from each other in the format of directory information and are not interchangeable with one another. Therefore, in a computer system using these standard OSes, when the optical card is read and written, access cannot be effected unless the data format in which data is recorded in the directory storage area 2b2 on the optical card coincides with the data format of the directory on the OS side.\nThat is, once the format of the directory storage area 2b2 on the optical card 1 is determined, operating systems which can access the optical card 1 are limited.\nGenerally, the optical card 11 is often carried with the user and used in various places. Consequently, he may encounter the case where his card must be read from or written into with an optical card recording/reproducing apparatus of a computer system running on an OS different from the OS used in recording the data on the optical card."}
-{"text": "The present invention relates to a discharge lamp apparatus that uses a discharge lamp as a light source, and particularly to an apparatus in which an electronic controller unit for applying a voltage to the discharge lamp is directly coupled with the discharge lamp.\nA discharge lamp apparatus that uses a discharge lamp as a light source is used as a vehicle headlight. An electronic controller unit that generates and controls a voltage applied to the discharge lamp includes a DC/DC converter for transforming an output voltage by switching an input voltage by a power device, a high voltage generation circuit for generating, from the output voltage of the DC/DC converter, a high voltage applied when lighting of the discharge lamp is initiated, and the like.\nWhen the temperature of electronic circuit components forming the controller unit rises, it is likely that, for instance, the soldered part of the circuit components melt and the circuit components operate erroneously. In the discharge lamp apparatus disclosed in JP-A-2000-235809, a part of a metallic heat radiator thermally coupled with a circuit substrate mounting circuit components thereon is exposed outside a headlight so that the heat generated by the discharge lamp and the circuit components may be radiated to the outside of the headlight through the metallic heat radiator.\nIn the discharge lamp apparatus disclosed in JP-A-2000-235809, an igniter part and a lighting device need be connected by a harness. This increases the number of component parts, complicates assembling work and adds manufacturing cost.\nIt is therefore proposed to directly couple and electrically connect the discharge lamp and the controller unit. However, if the discharge lamp and the controller unit are directly coupled and the controller unit is disposed near the discharge lamp, the internal temperature of the controller unit rises due to heat transferred or radiated from the discharge lamp and the heat generated by the controller unit itself. This is likely to cause erroneous operation of the circuit components in the controller unit.\nIt is therefore an object of the present invention to provide a discharge lamp apparatus that uses no high voltage wire nor high voltage connector and suppress rise of temperature of an electronic controller unit.\nIn a discharge lamp apparatus according to the present invention, a discharge lamp and an electronic controller unit for applying a voltage to the discharge lamp are directly coupled with and electrically connected to the discharge lamp. Therefore, a high voltage connector and a high voltage wire for connecting the discharge lamp and the controller unit are not necessitated.\nPreferably, a second casing mounting a power device of a DC/DC converter has a thermal conductivity higher than that of a first casing coupled with the discharge lamp. As a result, heat generated by the discharge lamp is less likely to be transferred from the first casing to the second casing, and the heat of the discharge lamp is less likely to be transferred circuit components mounted in the second casing.\nIn addition, heat generated by the power device of the DC/DC converter is more likely to be radiated to the outside of the second casing from the second casing that has the thermal conductivity higher than that of the first casing. Because the power device of the DC/DC converter generate more heat among the controller unit, the heat generated by the power device is readily radiated from the second casing to the outside of the second casing, thus suppressing rise of temperature of the circuit components in the controller unit including the power device. Thus, erroneous operation of the circuit components is prevented."}
-{"text": "In memory circuits utilizing variable threshold devices, such as MNOS transitors, the threshold value at which the transistor conducts may be controlled to have a low threshold level for supplying a logic or binary \"one\"; and to have a high threshold level fro supplying a logic or binary \"zero\", for example. In such circuits, it is desirable to have a short write cycle time to enable fast writing of data into the memory cell as well as reading such data out. In a prior art MNOS memory circuit, data may be written into the memory cell over and over. In other words, the memory cell may be written as a binary one for a number of times in a row before a binary zero is written. This causes the variable threshold transistor to shift its threshold between high and low to the maximum extent during the write cycle which saturates the device to a condition. which may be referred to as the saturated threshold state. Thus, when it is desired to write a binary level of the opposite state into the memory cell, the write pulse must have adequate polarizing voltage and time to shift the threshold voltage of the transistor from the previous saturated state. To speed up the time of writing a binary level or shifting the voltage in a variable threshold transistor, such as an MNOS transistor, the voltage or electric field across the gate insulator of the device is increased by increasing the polarization voltage. Unfortunately, these high electric fields across the gate insulator of the MNOS device during the write cycle serve to accelerate the undesirable endurance phenomena creating a smaller threshold voltage and decreased retention time of the memory state.\nIn an attempt to prevent saturation of the transistor, one typical prior art circuit operated such that the transistor was cleared prior to each writing of the transistor to a high threshold state. Such clearing consisted of applying a voltage to obtain a low threshold state prior to each writing of the data, regardless of the data already in the memory cell. This was followed by applying the voltage for the high threshold state to selected memory cells. Such a system, which required the application of a write voltage first for logical zeroes, and then for logical ones, required more than twice the minimum write time.\nIn order to overcome the disadvantages of operation created by the saturated threshold state without increasing the polarization voltage or without requiring that each cell be operated to its low threshold state prior to selectively operating the memory cells to the high threshold state during each write cycle, a system and method described in U.S. Pat. No. 4,090,258, dated May 16, 1978 and issued to J. R. Cricchi, one of the inventors of the present application, is proposed. This system describes a variable threshold MNOS non-volatile memory system and method that utilizes a pair of MNOS transistors for each memory cell. For a binary one, one transistor is operated to its high threshold state, while the other transistor is simultaneously operated to its low threshold state. For a binary zero, the one transistor is operated to its low threshold state while the other transistor is operated to its high threshold state. Prior to each write cycle of the circuit, a read-compare operation is performed, which compares the data to be written with the data actually stored in the memory cell. In the event that the data stored in the memory cell is the same as the data to be written into that memory cell, the write portion of the cycle is inhibited. This inhibiting of the rewriting of the same data into the memory cell prevents the transistor from being operated to its saturation threshold state without the necessity of clearing the transistor prior to the writing in of each new data. Further, in such system the fact that one of the transistors of the cell is operated to its low threshold state and the other to its high threshold state increases the \"window\" to minimize the errors when reading out the data stored in the memory cell.\nHowever, MNOS memories that have fast write characteristics tend to have short retention periods. Although the proposed system provides for a large window, it is possible that certain memory cell remain in their binary one or binary zero state during extended operations that the detection window would decrease; that is, the high threshold state of one of the transistors of the memory cell would decrease and the low threshold state of the other transistor of the memory cell would increase. Although such a system provided accuracy of operation down to a one-half volt difference between the two transistors of a memory cell, it is desirable that the transistors be refreshed periodically without driving the transistors to their saturated threshold state. Also, it is desirable that the refreshing of the memory cells be controlled internal to the memory array rather than by external programming."}
-{"text": "Internal combustion engines such as diesel engines, gasoline engines, and gaseous fuel-powered engines exhaust a complex mixture of air pollutants as byproducts of the combustion process. These air pollutants are composed of gaseous compounds including, among other things, the oxides of nitrogen (NOx). Due to increased attention on the environment, exhaust emission standards have become more stringent and the amount of NOx emitted to the atmosphere from an engine can be regulated depending on the type of engine, size of engine, and/or class of engine.\nOne method that has been utilized by engine manufacturers to comply with the regulation of engine emissions has been to implement exhaust gas regeneration (EGR). EGR systems recirculate exhaust gas to mix with combustion air supplied to an engine. The recirculated exhaust gas reduces a concentration of oxygen and increases a thermal mass within the engine's cylinders, thereby lowering a resulting combustion temperature. The lowered combustion temperature slows the chemical reaction of the combustion process and decreases the formation of NOx.\nAlthough successful at reducing NOx, the effectiveness of an EGR system can be influenced by an amount of heat contained within the exhaust gas that is recirculated. That is, the exhaust gas that is recirculated back into the engine can contain large amounts of energy in the form of waste heat, which can increase the in-cylinder combustion temperature and counteract the desired NOx reduction. In addition, the excess heat also taxes the engine's cooling system and may represent an unutilized source of energy that, if harnessed, could increase engine efficiency.\nAn exemplary system implementing EGR is disclosed in U.S. Pat. No. 6,216,458 of Alger et al. that issued on Apr. 17, 2001 (\u201cthe '458 patent\u201d). Specifically, the '458 patent discloses an exhaust gas recirculating (EGR) system for a turbocharged internal combustion engine. The EGR system includes a recirculation conduit for recirculating a volume of exhaust gas from the exhaust manifold to the intake manifold, a cooling air conduit for transporting a flow of cooling air, and a gas-to-air heat exchanger disposed in operative association with the recirculation conduit and the cooling air conduit. The heat exchanger is adapted for cooling the volume of the recirculated exhaust gas in the recirculation conduit."}
-{"text": "Herbal extracts and formulations have been suggested for the trichological treatment of humans, particularly males, but not exclusively so. Efficacy varies from nil to some reduction in hair loss and/or noticeable hair regrowth on the heads of patients. However, the interaction of the constituent compounds in such formulations with one another is little understood. The respective modes of action of such formulations are quite diverse. Generally, such treatments are either topical or systemic.\nThe term alopecia is used in the art generally to refer to all conditions of baldness, including partial baldness, and the term alopecia is used in this specification to have the same meaning, as well as to refer to hair loss generally, unless the context requires otherwise. The term \u201ctrichological treatment\u201d is used to refer to the treatment and/or formulations for the treatment of alopecia and hair thinning in males and females. Although the term is used sometimes to refer to scalp conditions or the like, the invention is not limited to treatment of such conditions, and may be used to treat hair loss on other parts of the body.\nThere are several root causes of alopecia, some of which are well established both clinically and pathologically, and quite well understood. Other root causes are not well understood and/or subject of conjecture. Some herbal and drug treatments of alopecia are directed to one or more of the root causes of alopecia, whether they be well understood or merely conjectured. There is a quantity of research data on efficacy, some of which is in conflict, indicating that the alopecia condition is complex, and is probably affected by a range of different mechanisms. Some root causes are expressed or exhibited more than others, but unless several of such root causes are addressed by a proposed remedy, it is suggested that its efficacy may be limited.\nThe present invention aims to provide formulations and treatments for trichology which alleviate one or more of the shortcomings of the prior art. Other aims and advantages of the invention may become apparent from the following description."}
-{"text": "1. Field of the Invention\nThe present invention generally relates to generating satellite tracking information for earth orbiting satellites. More specifically, the invention relates to a method and apparatus for generating satellite tracking information in a first format (e.g., a compact ephemeris model) through a network or communications link, then representing the satellite tracking information in a second format (e.g., a standard ephemeris model) at a receiver.\n2. Description of the Related Art\nA positioning receiver for the Global Positioning System (GPS) uses measurements from several satellites to compute a position. The process of acquiring the GPS radio signal is enhanced in speed and sensitivity if the GPS receiver has prior access to a model of the satellite orbit and clock. This model is broadcast by the GPS satellites and is known as ephemeris or ephemeris information. Each satellite broadcasts its own ephemeris once every 30 seconds. Once the GPS radio signal has been acquired, the process of computing position requires the use of the ephemeris information.\nThe broadcast ephemeris information is encoded in a 900 bit message within the GPS satellite signal. It is transmitted at a rate of 50 bits per second, taking 18 seconds in all for a complete ephemeris transmission. The broadcast ephemeris information is typically valid for 2 to 4 hours into the future (from the time of broadcast). Before the end of the period of validity the GPS receiver must obtain a fresh broadcast ephemeris to continue operating correctly and produce an accurate position. It is always slow (no faster than 18 seconds), frequently difficult, and sometimes impossible (in environments with very low signal strengths), for a GPS receiver to download an ephemeris from a satellite. For these reasons it has long been known that it is advantageous to send the ephemeris to a GPS receiver by some other means in lieu of awaiting the transmission from the satellite. U.S. Pat. No. 4,445,118, issued Apr. 24, 1984, describes a technique that collects satellite orbit information at a GPS reference station, and transmits the information to the remote GPS receiver via a wireless transmission. This technique of providing the ephemeris, or equivalent data, to a GPS receiver has become known as xe2x80x9cAssisted-GPSxe2x80x9d. Since the source of ephemeris in Assisted-GPS is the satellite signal, the ephemeris information remains valid for only a few hours. As such, the remote GPS receiver must periodically connect to a source of ephemeris information whether that information is received directly from the satellite or from a wireless transmission. Without such a periodic update, the remote GPS receiver will not accurately determine position.\nFurthermore, the Assisted-GPS systems typically retransmit the entire ephemeris message to the remote receiver. In many instances, bandwidth or packet size for the transmission of this message is not readily available.\nTherefore, there is a need for a method and apparatus for providing satellite trajectory and clock information to a remote receiver in a compact form.\nThe present invention is a method and apparatus for generating satellite tracking data (STD), then transmitting the data to a remote receiver in a compact form. The STD is derived by receiving at one or more satellite tracking stations the signals from at least one satellite and determining satellite tracking information (STI) through signal processing or by extracting the ephemeris message from the received signals. STI contains present satellite orbit trajectory data and satellite clock information.\nThe STD is reformatted into a compact format and provided to a remote satellite signal receiver via a network or communications system. The receiver converts the compact format into a standard format and uses the STD to compute the position of the receiver. The satellite system may include the global positioning system (GPS), GLONASS, GALILEO, or other satellite systems that may use STD to enhance the performance of the receiver."}
-{"text": "Starches from grain, cereals, and tubers, e.g., cornstarch, are widely used in the industrial manufacture of products such as sugar syrups and biofuels. For example, high fructose corn syrup (HFCS) is a processed form of corn glucose syrup having high fructose content and a sweetness comparable to sucrose, making HFCS useful as a sugar substitute in soft drinks and other processed foods. HFCS production currently represents a billion dollar industry. Similarly, the production of ethanol from starches is a rapidly expanding industry.\nSyrups and biofuels can be produced from starch by an enzymatic process that catalyzes the breakdown of starch into glucose. This enzymatic process typically involves a sequence of enzyme-catalyzed reactions:\n(1) Liquefaction: Alpha-amylases (EC 3.2.1.1) first catalyze the degradation of a starch suspension, which may contain 30-40% w/w dry solids (ds), to maltodextrans.\nAlpha-amylases are endohydrolases that catalyze the random cleavage of internal \u03b1-1,4-D-glucosidic bonds. Because liquefaction typically is conducted at high temperatures, e.g., 90-100\u00b0 C., thermostable alpha-amylases, such as alpha-amylases from Bacillus sp., are preferred for this step. Alpha-amylases currently used for this step, e.g., alpha-amylases from B. licheniformis, B. amyloliquefaciens, and Geobacillus stearothermophilus (AmyS), do not produce significant amounts of glucose. Instead, the resulting liquefact has a low dextrose equivalent (DE), containing maltose and sugars with high degrees of polymerization (DPn).\n(2) Saccharification: Glucoamylases catalyze the hydrolysis of alpha-1,4-glucosidic linkages of maltodextrins formed after liquefaction from non-reducing ends, releasing D-glucose. Saccharification produces high glucose syrup. Debranching enzymes, such as pullulanases, can aid saccharification.\n(3) Further Processing: A branch point in the process occurs after the production of a glucose-rich syrup. If the final desired product is a biofuel, yeast can ferment the glucose-rich syrup to ethanol. On the other hand, if the final desired product is a fructose-rich syrup, glucose isomerase can catalyze the conversion of the glucose-rich syrup to fructose.\n\nAlpha-amylases are isolated from a wide variety of bacterial, fungal, plant, and animal sources. Many industrially important alpha-amylases are isolated from Bacillus sp., in part because of the generally high capacity of Bacillus to secrete amylases into the growth medium. In addition, Bacillus alpha-amylase variants with altered while more desirable properties are obtained through genetic engineering. Furthermore, there is a need for blends of alpha-amylases, or variants thereof, which can capitalize on the best properties of at least two alpha-amylases of different origins.\nThe Fuelzyme\u00ae-LF alpha-amylase (SEQ ID NO: 2)(Verenium Corp.) is an engineered alpha-amylase obtained through DNA shuffling of three parental enzymes. See Richardson et al., J. Biol. Chem. 277: 26501-26507 (2002); U.S. Pat. No. 7,323,336. The advantageous properties of the Fuelzyme\u00ae-LF alpha-amylase include: effective viscosity reduction at a lower dose, improved thermostability, and broad pH operating ranges. The use of this alpha-amylase, however, is currently limited to biofuel applications, e.g., ethanol production, because it results in ineffectual glucose syrup that is not suitable for downstream applications such as sweetener applications. Specifically, saccharification of starch liquefact from Fuelzyme\u00ae-LF alpha-amylase results in iodine-positive saccharide (IPS), which indicates incomplete starch hydrolysis. Thus, if a way could be found to fully exploit the advantages of the Fuelzyme\u00ae-LF alpha-amylase in starch processing, particularly in sweetener applications, by using an optimized blend of alpha-amylases, this would also represent a useful contribution to the art."}
-{"text": "Cushioning materials have a variety of uses, such as for mattresses, seating surfaces, shoe inserts, packaging, medical devices, etc. Cushioning materials may be formulated and/or configured to reduce peak pressure on a cushioned body, which may increase comfort for humans or animals, and may protect objects from damage. Cushioning materials may be formed of materials that deflect or deform under load, such as polyethylene or polyurethane foams (e.g., convoluted foam), vinyl, rubber, springs, natural or synthetic fibers, fluid-filled flexible containers, etc. Different cushioning materials may have different responses to a given pressure, and some materials may be well suited to different applications. Cushioning materials may be used in combination with one another to achieve selected properties.\nFor example, cushioning materials may include a foam layer topped with a layer of thermoset elastomeric gel, such as a polyurethane gel or a silicone gel. Because polyurethane gels and silicone gels are generally structurally weak and/or sticky, cushioning materials may include film covering such gels, such as a thin thermoplastic polyurethane film. The film may reinforce the strength of the gel, and may prevent other materials from sticking to the gel, since the film generally adheres to the gel but is not itself sticky.\nGels may be used for cushioning and/or temperature management. Gels may provide cushioning because the gels may hydrostatically flow to the shape of a cushioned object and may tend to relieve pressure peaks. Gels may also reduce stresses from shear. Gels may have high thermal mass and/or thermal conductivity, and may therefore be used for heating (such as in hot packs for sore muscles), cooling (such as in cold packs for sprains or for a feeling of coolness when lying on a mattress or pillow), or maintaining a given temperature (such as in a mattress being used in a warm or cool room). For example, gel may be fused to the top of a mattress core, and a film may cover the gel. As another example, gels may be used as the top layer of a gel-on-foam wheelchair cushion.\nA conventional gel layer, with or without a plastic film, may be a barrier to gases (e.g., air, vapors, or other gases). This barrier may cause difficulties such as discomfort, such as when body heat and/or perspiration accumulate between the user's body and the gel layer. Even when a breathable material (such as a cover comprising foam or batting fiber) is disposed between a cushioned object and the gel, gases can only travel laterally through the breathable material. Since gases cannot penetrate the plastic film or the gel, the plastic film or the gel inhibits the flow of the gases away from the cushioned object. When the weight of the cushioned object compresses the breathable material, the lateral gas flow paths may become more constricted. Thus, it would be beneficial to provide a cushioning material that alleviates some of these concerns."}
-{"text": "The technology relates to an exposure device including a plurality of light-emitting elements and to a light receiving device including a plurality of light-receiving elements.\nAn exposure head including a plurality of light-emitting elements such as light-emitting diode (LED) elements is used in various image forming apparatuses. Non-limiting examples of such image forming apparatuses may include an electronic printer and a facsimile apparatus that each form an image by electrophotography scheme. For example, reference is made to Japanese Unexamined Patent Application Publication No. 2007-81081."}
-{"text": "The present invention relates to a dashboard support, and to a method for the production of a dashboard support.\nThe following discussion of related art is provided to assist the reader in understanding the advantages of the invention, and is not to be construed as an admission that this related art is prior art to this invention.\nModern motor vehicles have many driver assist systems as well as comfort and safety systems which are predominantly operated from the inside, in particular via the driver or passenger seats of a motor vehicle. Design features whether relating to the outer look or the interior space gain also increasingly on importance. For that reason, instrument panels, arranged in the motor vehicle between the front A-pillars beneath the windshield, are manufactured of plastic materials and sometimes covered with leather or other coating materials so as to receive a pleasing surface. To stiffen the instrument panel and to render it suitable for supporting various components, a dashboard support is arranged between the A-pillars. The dashboard support represents an elongate, rod-like structure having various mounting areas, for example, for steering column, airbags, glove compartment, heads-up display, navigation system and cable harness. Dashboard supports are typically coupled with the engine compartment rear bulk or the transmission tunnel for increase in stiffness.\nMotor vehicle parts are normally manufactured using manipulators or industrial robots to maneuver the various body components. The various parts are hereby assembled by a multi-stage fully automated process, requiring the parts to be aligned relative to one another for subsequent joining by welding or bonding for example. The body components have openings for engagement by grippers of the manipulators or industrial robots for handling purposes. Dashboard supports are normally handled in the area of the mounting areas for movement and installation in the region of the A-pillars.\nIt would be desirable and advantageous to provide an improved method for the production of a dashboard support to obviate prior art shortcomings."}
-{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and method of recording an image, and more particularly to an image recording apparatus and a method of recording an image which an image formed on a heat-development type photosensitive material by exposure is subjected to heat-development and is then transferred to an image receiving material so as to record the image thereon.\n2. Description of the Related Art\nIn a conventional image recording apparatus, an image formed on a heat-development type photosensitive material is subjected to heat-development and is then transferred to an image receiving material so as to record the image thereon. The image recording apparatus is provided with a photosensitive material magazine for storing a rolled heat-development type photosensitive material. The photosensitive material is pulled out from the photosensitive material magazine by a predetermined length necessary for forming an image thereon by exposure. The photosensitive material thus pulled out is cut and is then conveyed to an exposure section where an image is formed on the photosensitive material by exposure. After completion of the exposure, the photosensitive material is conveyed to a heat-development/transfer section via a water applying section. The image recording apparatus is also provided with an image receiving material magazine for storing a rolled image receiving material. The image receiving material is pulled out from the image receiving material magazine and is then conveyed to the heat-development/transfer section. In the heat-development/transfer section, the photosensitive material and image receiving material are conveyed around a cylindrical heating drum in a state in which they closely contact each other. The photosensitive material and image receiving material conveyed around the heating drum are both heated by heat from a halogen lamp built into the heating drum. As a result, the image formed on the photosensitive material by exposure is heat-developed and is transferred to the image receiving material.\nIn the heat-development section of the image recording apparatus, the photosensitive material and the image receiving material must be conveyed around the cylindrical heating drum. Therefore, a large space is needed to convey the photosensitive material and the image receiving material around the heating drum. Also, since the heating drum has a cylindrical shape, wasted space exists around the heating drum.\nAlso, the image receiving material magazine of the image recording apparatus must have a cylindrical space therein to store a rolled image receiving material. Accordingly, a large space is needed and wasted space exists therein.\nIn the image recording apparatus, the photosensitive material is pulled out from the photosensitive material magazine by a predetermined length necessary for forming an image thereon by exposure, cut, and then conveyed to the exposure section in which an image is formed on the photosensitive material by exposure. Consequently, a large space is necessary to pull out the photosensitive material from the magazine by a predetermined length necessary for forming an image thereon by exposure.\nFor the above-described reasons, the conventional image recording apparatus has the drawbacks of having a large size and high costs."}
-{"text": "1. Technical Field\nThe present disclosure is directed to a dynamic stabilization device and system for spinal implantation and, more particularly, to a dynamic stabilization device and system that is adapted to be positioned/mounted relative to first and second laterally-spaced pedicle screws and that includes at least one dynamic stabilization member that is positioned beyond the region defined between the pedicle screws, e.g., in an \u201coverhanging\u201d orientation.\n2. Background Art\nLow back pain is one of the most expensive diseases afflicting industrialized societies. With the exception of the common cold, it accounts for more doctor visits than any other ailment. The spectrum of low back pain is wide, ranging from periods of intense disabling pain which resolve, to varying degrees of chronic pain. The conservative treatments available for lower back pain include: cold packs, physical therapy, narcotics, steroids and chiropractic maneuvers. Once a patient has exhausted all conservative therapy, the surgical options range from micro discectomy, a relatively minor procedure to relieve pressure on the nerve root and spinal cord, to fusion, which takes away spinal motion at the level of pain.\nEach year, over 200,000 patients undergo lumbar fusion surgery in the United States. While fusion is effective about seventy percent of the time, there are consequences even to these successful procedures, including a reduced range of motion and an increased load transfer to adjacent levels of the spine, which accelerates degeneration at those levels. Further, a significant number of back-pain patients, estimated to exceed seven million in the U.S., simply endure chronic low-back pain, rather than risk procedures that may not be appropriate or effective in alleviating their symptoms.\nNew treatment modalities, collectively called motion preservation devices, are currently being developed to address these limitations. Some promising therapies are in the form of nucleus, disc or facet replacements. Other motion preservation devices provide dynamic internal stabilization of the injured and/or degenerated spine, without removing any spinal tissues. A major goal of this concept is the stabilization of the spine to prevent pain while preserving near normal spinal function. The primary difference in the two types of motion preservation devices is that replacement devices are utilized with the goal of replacing degenerated anatomical structures which facilitates motion while dynamic internal stabilization devices are utilized with the goal of stabilizing and controlling abnormal spinal motion.\nOver ten years ago a hypothesis of low back pain was presented in which the spinal system was conceptualized as consisting of the spinal column (vertebrae, discs and ligaments), the muscles surrounding the spinal column, and a neuromuscular control unit which helps stabilize the spine during various activities of daily living. Panjabi M M. \u201cThe stabilizing system of the spine. Part I. Function, dysfunction, adaptation, and enhancement.\u201d J Spinal Disord 5 (4): 383-389, 1992a. A corollary of this hypothesis was that strong spinal muscles are needed when a spine is injured or degenerated. This was especially true while standing in neutral posture. Panjabi M M. \u201cThe stabilizing system of the spine. Part II. Neutral zone and instability hypothesis.\u201d J Spinal Disord 5 (4):390-397, 1992b. In other words, a low-back patient needs to have sufficient well-coordinated muscle forces, strengthening and training the muscles where necessary, so they provide maximum protection while standing in neutral posture.\nDynamic stabilization (non-fusion) devices need certain functionality in order to assist the compromised (injured or degenerated with diminished mechanical integrity) spine of a back patient. Specifically, the devices must provide mechanical assistance to the compromised spine, especially in the neutral zone where it is needed most. The \u201cneutral zone\u201d refers to a region of low spinal stiffness or the toe-region of the Moment-Rotation curve of the spinal segment (see FIG. 1). Panjabi M M, Goel V K, Takata K. 1981 Volvo Award in Biomechanics. \u201cPhysiological Strains in Lumbar Spinal Ligaments, an in vitro Biomechanical Study.\u201d Spine 7 (3):192-203, 1982. The neutral zone is commonly defined as the central part of the range of motion around the neutral posture where the soft tissues of the spine and the facet joints provide least resistance to spinal motion. This concept is nicely visualized on a load-displacement or moment-rotation curve of an intact and injured spine as shown in FIG. 1. Notice that the curves are non-linear; that is, the spine mechanical properties change with the amount of angulations and/or rotation. If we consider curves on the positive and negative sides to represent spinal behavior in flexion and extension respectively, then the slope of the curve at each point represents spinal stiffness. As seen in FIG. 1, the neutral zone is the low stiffness region of the range of motion.\nExperiments have shown that after an injury of the spinal column or due to degeneration, neutral zones, as well as ranges of motion, increase (see FIG. 1). However, the neutral zone increases to a greater extent than does the range of motion, when described as a percentage of the corresponding intact values. This implies that the neutral zone is a better measure of spinal injury and instability than the range of motion. Clinical studies have also found that the range of motion increase does not correlate well with low back pain. Therefore, the unstable spine needs to be stabilized especially in the neutral zone. Dynamic internal stabilization devices must be flexible so as to move with the spine, thus allowing the disc, the facet joints, and the ligaments normal physiological motion and loads necessary for maintaining their nutritional well-being. The devices must also accommodate the different physical characteristics of individual patients and anatomies to achieve a desired posture for each individual patient.\nWith the foregoing in mind, those skilled in the art will understand that a need exists for a spinal stabilization device which overcomes the shortcoming of prior art devices. The present invention provides such an apparatus and method for spinal stabilization."}
-{"text": "1. Technical Field\nThe present invention relates to a technology that cleans the nozzle forming surface of a discharge head which discharges liquid, and, in particular, to a technology that performs wiping with a wiping member on the nozzle forming surface using cleaning liquid.\n2. Related Art\nIn the related art, a liquid discharging apparatus, such as an ink jet printer, which discharges liquid, such as ink or the like, from the nozzles of a discharge head is known. In such an apparatus, there is a case in which the liquid discharged from the nozzles adheres to the nozzle forming surface of the discharge head. If the liquid adheres to the nozzle forming surface, there is a problem in that the liquid is not appropriately discharged from the nozzles, thereby causing deterioration of image quality. Here, various technologies for removing liquid which adheres to the nozzle forming surface of a discharge head are proposed.\nFor example, a head cleaning mechanism disclosed in JP-A-2009-233896 is provided with an ejecting unit that ejects cleaning liquid onto a nozzle forming surface (droplet discharge surface) in order to suitably remove liquid which adheres to the nozzle forming surface. Further, after the cleaning liquid is ejected onto the nozzle forming surface, wiping is performed by a wiping member, and thus it is possible to suitably clean the nozzle forming surface.\nAs described above, in the head cleaning apparatus disclosed in JP-A-2009-233896, the cleaning liquid is directly ejected onto the nozzle forming surface, and thus the cleaning liquid easily comes in nozzles. Therefore, for example, there is a problem in that meniscus which are formed in ink in the vicinity of the openings of the nozzles are destroyed by the cleaning liquid which comes in the nozzles, and thus it is difficult to appropriately discharge ink from the nozzles. Here, in JP-A-2009-233896, it is described that the cleaning liquid is ejected onto the region of the nozzle forming surface other than the nozzles, and thus the cleaning liquid is prevented from coming in the nozzles. However, if the cleaning liquid is ejected onto the nozzle forming surface, there is limitation of preventing the cleaning liquid from coming in the nozzles."}
-{"text": "Vehicles are being provided with more complex systems. For example, vehicles commonly include a plethora of entertainment systems, such as stereos, USB interfaces for mobile telephones, video players, etc. Vehicles often have a host of other operator interfaces, such as emergency calling systems, vehicle navigation systems, heating and air conditioning systems, interior and exterior lighting controls, air bags, seatbelts, etc.\nVehicle operating environments are becoming more complex as well. For example, some roadways include u-turn lanes, round-a-bouts, no-left turn, multiple lanes one way in the morning and the other way in the afternoon, etc. Increases in traffic are also contributing to increased complexity.\nThese additional complexities contribute to increases in driver risk. What is needed are methods and systems for generating data representative of vehicle in-cabin insurance risk evaluations based on data representative of skeletal diagrams of a driver that are indicative of degrees of driver risk."}
-{"text": "The present invention relates to model aircraft and in particular to the remote radio control of such aircraft by means of series of PDM signals.\nOne of the major causes of accidents involving model airplanes is loss of control. Each year deadly accidents occur as a result of persons being struck by such models or slashed by the rotating propeller blades of a model out of control.\nTests conducted in Germany show that in 80% of model airplane accidents the aircraft remote control apparatus was intact after the accident indicating that the malfunction that lead to the accident was caused by interference with control signal transmission. The loss of communication between the model operator on the ground and his aircraft is thus believed to be a major cause of model aircraft accidents.\nIn view of the above, it is the principal object of the present invention to provide a system for the safe control of model aircraft in the event of a loss of communication between a model operator and his aircraft.\nA further object is to provide such a system which automatically shifts the model into a predetermined pattern in the event of a loss of communication which permits the model to be safely brought to the ground."}
-{"text": "The following includes information that may be useful in understanding the present invention(s). It is not an admission that any of the information provided herein is prior art, or material, to the presently described or claimed inventions, or that any publication or document that is specifically or implicitly referenced is prior art.\n1. Field of the Invention\nThe present invention relates generally to the field of support structures for safely carrying weapons and more specifically relates to a gun mounting apparatus for utility terrain vehicles.\n2. Description of the Related Art\nThe use of ATVs and UTVs has become increasingly popular. Given that weapons are often used in outdoor activities, many people install gun racks or gun mounts to their ATVs or UTVs to safely transport their weapons. Weapon racks and mounts for vehicles' prime function is to safely secure and transport weapons. The use of weapon racks and mounts for vehicles may further protect weapons against damage and accidental discharge which may occur if a weapon is not secured during transportation. Weapon racks and mounts should also facilitate uncomplicated, rapid, and hassle free access to the weapon.\nVarious gun racks and mounts have been proposed for vehicles. Unfortunately, many traditional gun racks that are currently available may be cumbersome and unsafe. Many may take up space in the personal cabin of vehicles, causing passengers to shift uncomfortably to avoid the guns. This also threatens the safety of passengers, as the guns or racks may strike them while vehicles are in motion. Many gun racks may be difficult to install, and may not allow for efficient application and removal of guns. Many gun racks and mounts comprise many parts, making it complicated to install or uninstall them from vehicles and to store them.\nTherefore a need exists for a gun mount for a vehicle which comprises few parts, is easy to install, keeps the weapon in a safe and comfortable position for driver and passenger and which simultaneously provides for easy access to the weapon and protects the weapon from accidental discharge and from being scratched or damaged.\nVarious attempts have been made to solve the above-mentioned problems such as those found in U.S. Pat. Nos. 5,495,969, 3,876,078, 2,251,271, 6,986,446, 7,137,511, and 4,995,537. This prior art is representative of support structures for safely carrying weapons. None of the above inventions and patents, taken either singly or in combination, is seen to describe the invention as claimed.\nIdeally, a gun mounting apparatus for utility terrain vehicles should be easy to assemble, comprise few parts, be easy to install and, yet, operates reliably and be manufactured at a modest expense. Thus, a need exists for a reliable UTV Gun Mount System to securely store and transport weapons and to avoid the above-mentioned problems."}
-{"text": "Crude oil produced from geological formations can contain various amounts of water. Water and crude oil are naturally non-miscible; however, when naturally occurring interfacial active compounds are present, these compounds can aggregate on the water and oil interface and cause water to form droplets within the bulk oil phase. During crude oil lifting through production tubings, the water and oil encounters an increased mixing energy from rapid flow through chokes and bends. This additional mixing energy can emulsify the water and oil. Such an oil external, water internal two-phase system is commonly referred to as a crude oil emulsion, which can be quite stable. The presence of water in crude oil, however, can interfere with refining operations, induce corrosion, increase heat capacity, and result in reduced handling capacity of pipelines and refining equipment. Therefore, the crude oil that is to be shipped out of the oilfield should be practically free of water and usually has a maximum water content limit of about 0.5 to 3% by total weight, depending on the type of crude and oil company.\nThe emulsified water can also contain various amounts of salts. These salts are detrimental to crude oil refining processes due to potential corrosion in the refinery. In crude oil refining, desalting techniques comprise the deliberate mixing of the incoming crude oil with a fresh \u201cwash water\u201d to extract the water soluble salts and hydrophilic solids from the crude oil. Primary dehydration of the crude oil occurs in oil field water oil separation systems such as \u201cfree water knock out\u201d and \u201cphase separators.\u201d Quite often, these systems are not adequate for efficient separation due to factors such as over production, unexpected production changes, and system underdesigns. In these cases, emulsion-breaking chemicals are added to the production processes to assist and promote rapid water oil separations.\nCommonly used emulsion-breaking chemicals or demulsifiers include alkylphenol formaldehyde resin alkoxylates (AFRA), polyalkylene glycols (PAG), organic sulfonates, and the like. These active ingredients are typically viscous and require a suitable organic solvent to reduce the viscosity of the demulsifier blend. Accordingly, there is an ongoing need for new, economical, environmentally-friendly, and effective chemicals and processes for resolving emulsions into the component parts of water and oil or brine.\nFor these reasons it is desired to have a demulsifier that does not require an organic solvent. It is also desired to have a demulsifier composition capable of resolving water external, water internal, and complex emulsions while having a broad dosage range."}
-{"text": "All publications herein are incorporated by reference to the same extent as if each individual publication or patent application was specifically and individually indicated to be incorporated by reference. The following description includes information that may be useful in understanding the present invention. It is not an admission that any of the information provided herein is prior art or relevant to the presently claimed invention, or that any publication specifically or implicitly referenced is prior art.\nMass spectrometry addresses two key questions: (1) \u201cwhat's in the sample?\u201d and (2) \u201chow much is there?\u201d. Both questions are addressed in the instant application. Several of the embodiments described herein focus on the first question; that is, identification of the components in a mixture. Embodiments of the present invention relate to software that has demonstrated substantial improvements in mass accuracy, sensitivity and mass resolving power. Certain of these gains follow directly from estimation and modeling of ion resonances using a physical model described by Marshall and Comisarow. Other embodiments described herein focus upon applications of estimation and modeling of the phases of ion resonances. Such methods can be divided into functional groups: phase-based methods, calibration, adaptive data-collection strategies, and miscellaneous auxiliary functions.\nThe traditional approach to analysis of Fourier transform mass spectrometry (\u201cFTMS\u201d) spectra is bottom-up. Resonances are detected in the spectra, from which inferences are made about the composition of the analyzed sample. Most of the embodiments described herein involve approaches to bottom-up analysis. Key steps in bottom-up analysis of FTMS data are detection and estimation of ion resonances, mass calibration, and identification. Various embodiments of the present invention involve reducing the 4 MB of data representing an FTMS (MS-1) spectrum to a list of candidate elemental compositions for each detected peak with probabilities assigned to these identities and abundance estimates. The essential information represents a data reduction of roughly three orders of magnitude relative to the unprocessed spectrum. In the bottom-up approach to data analysis, peaks are detected and characterized by estimation first, and then knowledge about the sample is used to calibrate and identify the components. The ability to perform these calculations in real-time creates exciting possibilities for adaptive workflows that actively direct acquisition of optimally informative data."}
-{"text": "1. Field of Invention\nThe present invention relates to fuel injection valves for use in internal combustion engines.\n2. Description of Related Art\nVarious types of fuel injection valves have been hitherto proposed and put into practical use particularly in the field of automotive internal combustion engines. One of them is disclosed in Laid-open Japanese Patent Application (Tokkai) 2000-8990, which generally comprises a valve element arranged in a casing to move between open and close positions and an electromagnetic coil arranged to actuate the valve element to move between the open and close positions in accordance with energization/de-energization of the electromagnetic coil. Upon opening of the valve element, a pressurized fuel in a fuel passage of the casing is injected into a given space, such as a combustion chamber, air intake passage or the like, of the internal combustion engine.\nHowever, due to inherent construction, the fuel injection valves of the type disclosed by the above-mentioned Japanese Application tend to fail to exhibit a satisfied dimensional stability of the valve element. That is, for producing the valve element, a press working is employed to shape a given portion of the valve element, which however tends to induce a deformation or swelling of the given portion to which the pressing force is actually applied. As is known, such deformation or swelling (viz., dimensional poorness) induces a non-smoothed movement of the valve element and thus tends to exhibit a poor responsiveness of the same upon energization/de-energization of the electromagnetic coil. If, for smoothing the movement, a finish machining working is additionally employed for finely shaping the valve element, productivity of the fuel injection valve is lowered and thus cost performance of the same becomes poor."}
-{"text": "Integrated circuit technology has revolutionized various fields, including computers, control systems, telecommunications, and imaging. One field in which integrated circuitry is widely used is video imaging. Different types of semiconductor imagers include: charge coupled devices, photodiode arrays, charge injection devices, and hybrid focal plane arrays. Many of these devices include pixels that are arranged in sensor arrays to convert light images into electrical signals.\nExamples of MOS imaging devices are detailed in \u201cA \u00bc Inch Format 250K Pixel Amplified MOS Image Sensor Using CMOS Process\u201d by Kawashima et al., IEDM 93\u2013575 (1993), and \u201cA Low Noise Line-Amplified MOS Imaging Devices\u201d by Ozaki et al., IEEE Transactions on Electron Devices, Vol. 38, No. 5, May 1991. In addition, U.S. Pat. No. 5,345,266 to Denyer, titled \u201cMatrix Array Image Sensor Chip,\u201d describes a MOS image sensor. The devices disclosed in these publications provide a general design approach to MOS imaging devices. In addition, MOS approaches to color imaging devices are described in \u201cColor Filters and Processing Alternatives for One-Chip Cameras,\u201d by Paruiski, IEEE Transactions on Electron Devices, Vol. ED-32, No. 8, August 1985, and \u201cSingle-Chip Color Cameras With Reduced Aliasing\u201d by Imaide et al., Journal of Imaging Technology, Vol. 12, No. 5, October 1986, pp. 258\u2013260.\nImage sensor circuitry generally includes circuits for performing black level calibration and automatic gain control. Black level calibration attempts to eliminate the portion of the image signal that exists when no light is being detected, thus allowing for a truer zero reference signal so that the later signal processing is improved. After the black level calibration has been performed, automatic gain control amplifies the video signal at a controlled level so as to utilize more of the available signal amplification range. The combination of the black level calibration and automatic gain control can be said to form one type of an \u201cauto black expansion\u201d method. One prior art circuit that addresses issues related to these processes is shown in U.S. Pat. No. 4,187,519 to Vitols et al. FIGS. 1, 2, 3, 4, 5, 6 and 7 of Vitols et al. have been reproduced herein as FIGS. 1, 2A, 2B, 2C, 2D, 2E, and 2F, respectively.\nVitols et al. disclose a circuit that is designed to provide a video contrast expansion system. As illustrated in FIG. 1, a sequence of video pixel data representative of an image is developed by a signal source 11, which may be, for example, a radar, an IR (infrared) sensor or a TV (television) camera. This pixel data is applied to an I-getter (intensity getter) 13 and to a delay circuit 19. An initialize command signal from the source 11 initializes the operation of the I-getter 13 before or during the time that the sequence of pixels is being generated by the source 11.\nThe I-getter 13 searches the input sequence of pixels (designated as new data or ND) to determine bias and gain intensity parameters. A bias function generator 15 is automatically adjusted by the bias parameters to develop pixel bias correction signals. In a similar manner a gain function generator 17, which is similar in structure and operation to the generator 15, is automatically adjusted by the gain parameters to develop pixel correction signals.\nIn its operation the I-getter 13 reduces and transforms the input sequence of pixels into a reduced number of bias parameters and a reduced number of gain parameters. Each of the function generators 15 and 17 smoothly fills in or shades in its associated widely spaced apart parameters (bias or gain) by double linear interpolations to produce its associated correction signals (bias or gain). Each sequence of correction signals is equal to the number of pixels in the input sequence of pixels.\nIt takes time for each of the function generators 15 and 17 to get sufficient parameter data in before it can start computing its associated correction signals. Furthermore, it takes additional time before each of the function generators 15 and 17 compute its associated correction signals. This combined delay time is offset by the delay circuit 19 which synchronizes the time of occurrence of the sequence of pixels at the output of the delay circuit 19 with the time of generation of the pixel bias correction signals and pixel gain correction signals.\nThe computed pixel bias correction signals are respectively subtracted in a combiner or subtractor 21 from the delayed or synchronized sequence of pixels to selectively lower the minimum values in associated groups of pixels in the delayed sequence of pixels to very close to zero. The minimized output of the combiner 21 is respectively multiplied in a multiplier or AGC (automatic gain control) circuit 23 to selectively expand the maximum amplitudes in the associated groups of pixels at the output of the combiner 21 to near the saturation level of the electronics of the system of FIG. 1. The output of the AGC circuit 23 can be applied to a display generator (not shown) to generate an enhanced picture of the image in which the video contrast of the image has been substantially expanded in two directions (minimum and maximum values of contrast).\nThe pixel data at the output of the signal source 11 may have a very narrow dynamic range of contrast, making it extremely difficult to discern objects in areas of lower contrast if directly viewed at the point. However, by the operation just explained, the pixel data at the output of the AGC circuit 23 has its dynamic range of contrast selectively expanded or stretched from near minimum to near maximum, making objects originally in low areas of contrast now clearly discernible when displayed.\nFIGS. 2A\u20132F illustrate some of the possible operations that the circuit of FIG. 1 can perform on a one-dimensional signal. Such operations are analogous to those performed on a two-dimensional signal since a two-dimensional signal is essentially comprised of a vertical plurality of horizontal one-dimensional signals.\nFIG. 2A illustrates an exemplary one-dimensional signal I(x) comprised of a sequence of 256 pixels. Although I(x) is shown as an analog signal it is a sequence of 256 intensity (I) points of varying amplitudes. Note that I(x) varies over a fairly wide intensity range from an amplitude of 630 (Isat or the intensity saturation level of the system) to an amplitude of about 50 at the low end.\nIt may be desired in some applications to look for small, rather than large, amplitude deviations from nominal. For example, if the intensity variations in the incoming signals I(x) of FIG. 2 represent pixels or intensity points of a scene, very high intensity portions of I(x) would ultimately develop very bright portions in a picture, and very low intensity portions of I(x) would ultimately develop very dark portions of the picture. However, a human observer or operator may not necessarily be interested in very high or very low intensity levels in the incoming signal I(x). Rather, an observer may only be interested in small variations of intensity in I(x), whether contained in very high, very low and/or intermediate levels of intensity in I(x). This is due to the fact that the small variations of intensity in I(x) can define the details of a scene or image sufficiently to possibly enable an observer to identify what is happening or contained in the picture.\nBy using a conventional contrast enhancement technique, an observer may be unable to discern what is contained in picture areas of small intensity variations. It is to the correction of this problem that the circuit of FIG. 1 is directed. To illustrate, assume that the low intensity region of pixels between 100 and 200 in FIG. 2 represents the signal that is coming from the inside of a dark cave. The brightest pixel in the 100\u2013200 region of pixels is dark in comparison with the other high intensity or bright pixels in FIG. 2. An observer may be looking for the glint of light on a gun barrel inside that cave. If the dynamics of the signal in that 100\u2013200 pixel region remain unchanged, an observer would not be able to see that glint of light on that gun barrel. However, if the gain of the signal in the 100\u2013200 pixel region (where the amplitude of the signal is very low) could be increased so that variations of the signal show up to substantially the fullest extent possible without saturating the electronics of FIG. 1, the human observer could more readily determine what was inside the cave. It should be noted that a resultant picture would not look like a normal picture since the inside of the cave would be just as bright as what surrounds the cave on the outside. However, the intent of the circuit of FIG. 1 is to brighten everything to the fullest dynamic range, so that an observer can see what is happening or contained in a resultant picture.\nFIG. 2B illustrates a possible guide as to what may be done to increase the dynamic range of I(x) of FIG. 2A. As shown, FIG. 2B illustrates a plurality of eight local data regions, each 32 pixels in length, which respectively encompass contiguous 32-pixel long portions of I(x). In each data region the brightest (or largest intensity) and darkest (or smallest intensity) parts of the associated portion of I(x) are determined. The largest and smallest intensity values found in each data region are used to form horizontal ceiling (C) and floor (F) values for each pixel across that region. Thus, the sequence of ceiling values found in the respective data regions forms a ceiling function of X values, designated C(x), while the sequence of floor values found in the respective data regions forms a floor function of X values, designated F(x), for the included portions of I(x). FIG. 2C shows just F(x) the floor function of FIG. 2B. As shown, F(x) is comprised of the segmented sequence of minimum values across the respective regions of FIG. 2B.\nAs described, a first possible operation that can be performed to maximize the dynamic range of I(x) is to subtract F(x) from I(x) to produce the wave form shown in FIG. 2D. It can be seen in FIG. 2D that the amplitude of the signal I(x) of FIG. 2A is substantially reduced, while still substantially retaining the intensity variations of FIG. 2A.\nAs also described, the next possible operation that can be performed is to expand the amplitude of the signal I(x)\u2013F(x) of FIG. 2D by adjusting the gain in the first data region (0\u201332), second data region (32\u201364), third data region (64\u201396) etc.\u2014each gain adjustment being independent of the others-so that the maximum value of the signal in each data region go all the way up to the maximum permissible level or saturation level (Isat) FIG. 2E illustrates a piecewise constant gain function G(x) that would independently adjust the amplitude in each segment or data region up to the saturation level. This piecewise constant gain function would be determined by the value of the saturation intensity (Isat) divided by the difference between the C(x) and F(x) functions.\nFIG. 2F illustrates the result of multiplying the function of FIG. 2D by the gain function of FIG. 2E. This multiplication raises the maximum value in each data region up to the saturation level. However, Vitols et al. describe FIG. 2F as being very unsatisfactory because of the discontinuities or, places where the amplitude of the signal rises or falls vertically. Although the full dynamic range of the signal is obtained in each data region, the resultant signal shown in FIG. 2F is very bumpy. So the overall information that an operator may be looking for in the signal may be totally lost in the discontinuities. If the one-dimensional wave form of FIG. 2F were applied to a two-dimensional picture, a very strong checkerboard pattern would result, which would distort the operator's perception of a picture or image to the point where he probably could not discern what he was looking at. Vitols et al. go on to explain how the discontinuities shown in FIG. 2F can be eliminated by smoothing the signal functions F(x) and G(x) from one data region to another. This smoothing operation is described further in the Vitols et al. patent, but will not be discussed further herein.\nOne of the drawbacks of the Vitols et al. circuit is that it requires a complex sequence of operations to determine the desired adjustment functions for the system. As described, the I-getter (as illustrated in FIG. 12 of Vitols et al.) requires two shift registers for every data cell used in an output picture. In addition, a large number of switching points capable of switching out each of the data cells for separate comparator operations must be used. In addition, during the comparator operations the comparator voltage level to which the incoming pixel signals are compared may be constantly shifting. In addition, as further described in Vitols et al., it takes time for each of the function generators to get sufficient parameter data in before the function generators can start computing the associated correction signals. Furthermore, it is also stated that it takes additional time before each of the function generators can compute its associated correction signals. This combined delay time must be offset by a delay circuit which synchronizes the time of occurrence of the sequence of pixels with the time of generation of the pixel bias correction signals and pixel gain correction signals.\nThe present invention is directed to a circuit that overcomes the foregoing and other problems in the prior art. More specifically, the present invention is directed to an auto black expansion method and apparatus for an image sensor that uses a simplified digital control system and does not require additional shift registers for the pixel signals or continual shifting of input comparator levels during a given field."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to a semiconductor memory device and a method of manufacturing the same. More particularly, the invention relates to a semiconductor memory device having an etch stopper formed in a cell region and a method of manufacturing the semiconductor memory device.\nA claim of priority is made to Korean Patent Application No. 10-2005-0015371, filed on Feb. 24, 2005, the disclosure of which is hereby incorporated by reference in its entirety.\n2. Description of Related Art\nAs the integration density of semiconductor memory devices increases, contact holes formed in the semiconductor devices become increasingly narrow. As a result, the aspect ratio of the contact holes tends to increase, and therefore contact holes and associated contacts are often required to be formed using multi-step processes.\nThe integration density of dynamic random access memory (DRAM) devices is typically very high. Because of this, active regions between adjacent gate electrodes in DRAMs are generally so narrow that it is difficult to form contacts connected to the active regions. One approach to forming contacts in narrow regions is to form a hard mask layer and spacers over and around the gate electrodes before forming self aligned contact pads between the gate electrodes. Then, bitline contacts or storage node contacts are formed over the self-aligned contact pads. This approach is often referred to as \u201cself aligned contact technology.\u201d\nTypically, the self-aligned contacts have a larger diameter than the bitline contacts or the storage node contacts. However, as the integration density of DRAM devices continues to increase, the diameter of the self-aligned contact pads reduces and approaches the diameter of the bitline and storage node contacts. As a result, slight misalignments between the bitline or storage node contacts and node can cause serious problems in the manufacture of DRAM devices. Illustrations of some common problems are described below in relation to FIG. 1.\nReferring to FIG. 1, a conventional semiconductor memory device comprises a semiconductor substrate 10 in which a cell region \u201cA\u201d and a core and peripheral region \u201cB\u201d are defined. A plurality of isolation layers 15 defining active regions 10a are formed in semiconductor substrate 10, and junction regions 25a and 25b are formed between isolation layers 15 in respective cell and core and peripheral regions \u201cA\u201d and \u201cB\u201d.\nA first insulating interlayer 30 is formed over semiconductor substrate 10 and a second insulating interlayer 40 is formed over first insulating interlayer 30.\nIn cell region \u201cA\u201d, a plurality of self-aligned contact pads 35 are formed in first insulating interlayer 30 over junction regions 25a, and a plurality of bitline contact holes 45 are formed in second insulating interlayer 40 to expose self-aligned contact pads 35.\nIn core and peripheral region \u201cB\u201d, a gate electrode structure 20 is formed in first insulating layer 30 on side junction region 25b and an interconnection contact hole 50 is formed in first and second insulating interlayer 30 and 40 to expose side junction region 25b. Each gate electrode structure 20 comprises a gate insulating layer 16 formed on semiconductor substrate 10, a gate electrode material 17 formed on gate insulating layer 16, a hard mask layer 18 formed over gate electrode material 17, and spacers 19 formed on the sidewall of gate insulating layer 16, gate electrode material 17, and hard mask layer 18.\nTypically, the respective diameters of self-aligned contact pads 35 and bitline contact holes 45 are slightly different from each other, or they are slightly misaligned with each other. As a result, portions of first insulating interlayer 30 adjacent to self-aligned contact pads 35 may be removed when bitline contact holes 45 are formed. Due to the partial loss of first insulating interlayer 30, active regions 10a may also be partially removed, causing leakage current in the semiconductor memory device.\nStorage node contact holes (not shown) are typically formed after bitlines are formed in bitline contact holes 45. The storage node contact holes are generally formed between the bitlines to contact self-aligned contact pads 35. However, the storage node contact holes are almost inevitably misaligned with the self-aligned contact pads. As a result, the first insulating interlayer 30 is almost inevitably etched when the storage node contact holes are formed, causing further damage to active regions 10a. \nInterconnection contact hole 50 is typically formed in core and peripheral region \u201cB\u201d at the same time when bit line contact holes 45 are formed in cell region \u201cA\u201d. Unfortunately, however, as the integration density of core and peripheral region \u201cB\u201d increases, interconnection contact hole 50 tends to expose more than just junction region 25b. For instance, interconnection contact hole 50 may expose parts of gate electrode structure 20.\nConsequently, spacers 19 and hard mask layer 18 may be partially lost during the formation of interconnection contact hole 50. As a result, gate electrode material 17 may be exposed and may be short-circuited by a conductive material formed later in interconnection contact hole 50.\nTo solve these problems, some conventional memory devices interpose etch stoppers between first insulating layer 30 and second insulating interlayer 40 to prevent first insulating interlayer 30 from being etched when bitline contact holes 45 are formed. A conventional memory device employing this technique is disclosed, for example, in U.S. Pat. No. 6,787,906 (the '906 patent).\nIn the '906 patent, etch stoppers are formed on first insulating layer 30. Accordingly, bitline contact holes 45 and storage node contact holes can be formed without further etching first insulating layer 30. However, the etch stoppers are also formed in core and peripheral region \u201cB\u201d, making it difficult to form interconnection contact hole 50. Since interconnection contact hole 50 penetrates first and second insulating layers 30 and 40, if the etch stoppers are interposed between first and second insulating interlayers 30 and 40, it is difficult to expose junction region 25b. "}
-{"text": "1. Field of the Invention\nThe present invention relates to a plasma display panel and the manufacturing method thereof, more particularly to the partition wall structure of the panel and the manufacturing method thereof.\n2. Description of the Prior Art\nThe rib of the plasma display panel (referred to PDP in the following) commonly has a stripe-shaped structure. However, the grid-mesh rib structure is also used at present, for example, the one disclosed in the U.S. Pat. No. 5,701,056 by NEC. The structure disclosed by NEC forms stripe-shaped ribs on the back substrate of the PDP and forms grid-mesh-shaped ribs on the front substrate of the PDP, then assembles the front and back substrates, as shown in FIG. 1. The structure disclosed by NEC has the following four disadvantages:\nSince the front substrate has an additional rib manufacturing process in the NEC structure, the cost is relatively high.\nWhen assembling the front and the back substrates, the high aligning precision of the two substrates is strictly required; this deepens the difficulty of the manufacturing process.\nTo ensure that the front and the back substrate are precisely aligned, increasing the width of the rib of the front and the back substrates is often required. Hence the opening rate of the PDP is compromised.\nDue to the width of the rib, the effective area of the coating fluorescent body becomes smaller."}
-{"text": "The present invention relates to a novel hard, highly abrasion and corrosion-resistant, hydrogenated carbon material useful as an ultra-thin protective overcoat layer for high areal recording density thin film magnetic and magneto-optical (xe2x80x9cMOxe2x80x9d) recording media, a method of depositing films or layers of the novel material, and to improved magnetic and MO media including a protective overcoat layer comprised of the novel material.\nA magnetic recording medium, e.g., a hard disk, typically comprises a laminate of several layers, comprising a non-magnetic substrate, such as of Alxe2x80x94Mg alloy or a glass or glass-ceramic composite material, and formed sequentially on each side thereof, a polycrystalline underlayer, typically of chromium (Cr) or Cr-baed alloy, a polycrystalline magnetic recording medium layer, e.g., of a cobalt (Co)-based alloy, a hard, abrasion-resistant, protective overcoat layer, typically containing carbon (C), and a lubricant topcoat.\nIn operation of the magnetic recording medium, the polycrystalline magnetic recording medium layer is locally magnetized by a write transducer, or write head, to record and store information. The write transducer creates a highly concentrated magnetic field which alternates direction based on the bits of information being stored. When the local magnetic field produced by the write transducer is greater than the coercivity of the recording medium layer, then the grains of the polycrystalline recording medium at that location are magnetized. The grains retain their magnetization after the magnetic field produced by the write transducer is removed. The direction of the magnetization matches the direction of the applied magnetic field. The magnetization of the polycrystalline recording medium can subsequently produce an electrical response in a read transducer, allowing the stored information to be read.\nThin film magnetic recording media are conventionally employed in disk form for use with disk drives for storing large amounts of data in magnetizable form. Typically, one or more disks are rotated on a central axis in combination with data transducer heads. In operation, a typical contact start/stop (CSS) method commences when the head begins to slide against the surface of the disk as the disk begins to rotate. Upon reaching a predetermined high rotational speed, the head floats in air at a predetermined distance from the surface of the disk due to dynamic pressure effects caused by air flow generated between the sliding surface of the head and the disk. During reading and recording operations, the transducer head is maintained at a controlled distance from the recording surface, supported on a bearing of air as the disk rotates, such that the head can be freely moved in both the circumferential and radial directions, allowing data to be recorded on and retrieved from the disk at a desired position. Upon terminating operation of the disk drive, the rotational speed of the disk decreases and the head again begins to slide against the surface of the disk and eventually stops in contact with and pressing against the disk. Thus, the transducer head contacts the recording surface whenever the disk is stationary, accelerated from the static position, and during deceleration just prior to completely stopping. Each time the head and disk assembly is driven, the sliding surface of the head repeats the cyclic sequence consisting of stopping, sliding against the surface of the disk, floating in the air, sliding against the surface of the disk, and stopping.\nAs a consequence of the above-described cyclic CSS-type operation, the surface of the disk or medium surface wears off due to the sliding contact if it has insufficient abrasion resistance or lubrication quality, resulting in breakage or damage if the medium surface wears off to a great extent, whereby operation of the disk drive for performing reading and reproducing operations becomes impossible. The protective overcoat layer is formed on the surface of the polycrystalline magnetic recording medium layer so as to protect the latter from friction and like effects due to the above-described sliding action of the magnetic head. Abrasion-resistant, carbon (C)-containing protective coatings have been utilized for this purpose, and are typically formed by sputtering of a carbon target in an argon (Ar) atmosphere. Such amorphous carbon (a-C)-containing protective overcoat layers formed by sputtering have relatively strong graphitic-type bonding, and therefore exhibit a low coefficient of friction in atmospheres containing water (H2O) vapor, which characteristic is peculiar to graphite. However, the a-C layers produced in such manner have very low hardness as compared with many ceramic materials such as are employed as slider materials of thin film heads, and thus are likely to suffer from wear due to contact therewith.\nIn recent years, therefore, carbon-based protective overcoat layers having diamond-like hardness properties (i.e., HV of about 1,000-5,000 kg/mm2) have been developed, and films of diamond-like carbon (DLC) having a high percentage of diamond-type Cxe2x80x94C bonding have been utilized. Such DLC films exhibit a high degree of hardness due to their diamond-like sp3 bonding structure, and in addition, exhibit the excellent sliding properties characteristic of carbon, thus affording improved sliding resistance against sliders composed of high hardness materials. Such DLC films are generally obtained by DC or RF magnetron sputtering of a carbon target in a gas atmosphere comprising a mixture of Ar gas and a hydrocarbon gas, e.g., methane, or hydrogen gas. The thus-obtained films exhibit DLC properties when a fixed amount of hydrogen is incorporated therein. Incorporation of excessive amounts of hydrogen in the films leads to gradual softening, and thus the hydrogen content of the films must be carefully regulated.\nAmorphous, hydrogenated carbon films (referred to herein as a-C:H films) obtained by sputtering of carbon targets in an Ar+H2 gas mixture exhibiting diamond-like properties have also been developed for improving the tribological performance of disk drives; however, the electrical insulating properties of such type films lead to undesirable electrical charge build-up or accumulation during hard disk operation which can result in contamination, glide noise, etc. In order to solve this problem without sacrifice or diminution of the advantageous mechanical properties of such a-C:H films, attempts have been made to dope or otherwise incorporate nitrogen (N) atoms into the a-C:H films, in view of a substantial decrease in electrical resistivity and optical band gap (EBG) exhibited by such nitrogen-doped a-C:H films relative to undoped films.\nHowever, the continuous increase in areal recording density of magnetic recording media requires a commensurately lower flying height. Therefore, it would be advantageous to reduce the thickness of the carbon-based protective overcoat layer without adverse consequences. Conventional sputtered a-C:H materials are difficult to uniformly deposit and generally do not function satisfactorily at a thickness of about 30 xc3x85 or less. Specifically, conventional sputtered a-C:H films of about 30 xc3x85 thickness fail to provide adequate protection against corrosion of the underlying magnetic layer(s), particularly Co-containing ferromagnetic layers, when under environments of high temperature and humidity, and the resulting corrosion product(s) frequently are disadvantageously transferred to the transducer heads, often leading to failure of the disk drive.\nThe use of alternative deposition techniques for developing thinner and harder a-C:H layers having the requisite mechanical and tribological properties has been studied, such as chemical vapor deposition (CVD), ion beam deposition (IBD), and cathodic arc deposition (CAD) techniques. For example, the IBD method can be utilized for forming hydrogenated ion-beam carbon films (referred to herein as i-C:H films) that exhibit superior tribological performance at thicknesses below about 100 xc3x85. However, such films are insulating and, thus, suffer from the above-described drawback of electrical charge build-up during hard disk operation associated with sputtered a-C:H films.\nAccordingly, there exists a need for an improved hard, abrasion and corrosion-resistant material particularly suitable for use as an ultra-thin protective overcoat layer in high areal density magnetic recording media, and a method for manufacturing same, which method is simple, cost-effective, and fully compatible with the productivity and throughput requirements of automated manufacturing technology.\nThe present invention fully addresses and solves the above-described problems attendant upon the formation of ultra-thin, abrasion and corrosion-resistant protective overcoat layers suitable for use with high areal density magnetic recording media, such as are employed in hard drive applications, while maintaining full compatibility with all mechanical and electrical aspects of conventional disk drive technology. In addition, the present invention enjoys utility in the formation of ultra-thin, abrasion and corrosion-resistant protective overcoat layers required in the manufacture and use of thin film-based, ultra-high recording density magneto-optical (MO) data/information storage and retrieval media in disk form and utilizing conventional Winchester disk drive technology with laser/optical-based read/write transducers operating at flying heights on the order of a few micro-inches above the media surface.\nAn advantage of the present invention is an improved hard, abrasion and corrosion-resistant, hydrogenated carbon (xe2x80x9cC:Hxe2x80x9d) material formed by a simultaneous sputter+plasma enhanced chemical vapor deposition (xe2x80x9cPECVDxe2x80x9d) process.\nAnother advantage of the present invention is an improved protective overcoat material for magnetic and MO recording media and comprising an improved hard, abrasion and corrosion-resistant, sputter+PECVD-deposited C:H material.\nYet another advantage of the present invention is an improved magnetic or MO recording medium including an ultra-thin protective overcoat layer comprised of a sputter+PECVD-deposited C:H material.\nStill another advantage of the present invention is an improved simultaneous sputter+PECVD method for forming C:H films or layers suitable for use as ultra-thin, abrasion and corrosion-resistant protective overcoat materials in magnetic and MO recording media applications.\nA further advantage of the present invention is an improved method for regulating or controlling the amount of PECVD-derived carbon (C) atoms contained in sputter+PECVD deposited C:H materials, films or layers.\nA still further advantage of the present invention is an improved apparatus for performing simultaneous sputtering+PECVD for forming films or layers of the improved C:H materials according to the present invention.\nStill another advantage of the present invention is an improved magnetic recording medium comprising at least one ferromagnetic layer containing Co and means for protecting the at least one Co-containing ferromagnetic layer from corrosion under high temperature, high humidity conditions.\nAdditional advantages and other features of the present invention will be set forth in the description which follows and in part will become apparent to those having ordinary skill in the art upon examination of the following or may be learned from the practice of the present invention. The advantages of the present invention may be realized and obtained as particularly pointed out in the appended claims.\nAccording to one aspect of the present invention, the foregoing and other advantages are obtained in part by a novel hard, abrasion and corrosion-resistant material useful as an ultra-thin protective overcoat layer for a magnetic or magneto-optical (MO) recording medium, which novel material comprises hydrogenated carbon (C:H) formed by a process comprising simultaneous sputter and plasma-enhanced chemical vapor deposition (PECVD) of the hydrogenated carbon (C:H), wherein the amount of carbon atoms in the novel C:H material derived from the PECVD component of the process is less than about 50 at. %.\nAccording to certain embodiments of the present invention, the amount of carbon atoms in the novel C:H material derived from the PECVD component of the process is at least about 30 at. %; the position of the Raman G-band of the C:H material is about 1553 cmxe2x88x921; and the film resistance of the C:H material is as high as about 85 k xcexa9.\nAnother aspect of the present invention is a magnetic or MO recording medium comprising a protective overcoat layer formed of the novel C:H material; and according particular embodiments of the present invention, the thickness of the protective overcoat layer is not greater than about 30 xc3x85.\nYet another aspect of the present invention is a method of forming a layer of a novel hard, abrasion, and corrosion-resistant hydrogenated carbon (C:H) material on a surface of a substrate, which method comprises the steps of:\n(a) providing a vacuum chamber including a carbon sputtering target in the interior space thereof;\n(b) providing a substrate in the interior space of the chamber, such that a surface of the substrate is in facing relation to the sputtering target;\n(c) supplying the interior space of the vacuum chamber with at least one hydrocarbon gas and at least one inert gas at separately controllable flow rates and applying a sufficient negative potential to the carbon sputtering target to generate a plasma in the interior space to deposit a layer of a novel hard, abrasion and corrosion-resistant C:H material on the substrate surface by simultaneous sputtering of the carbon sputtering target and plasma enhanced chemical vapor deposition (PECVD) of carbon and hydrogen from the at least one hydrocarbon gas, wherein:\nstep (c) includes separately controlling the flow rates of each of the hydrocarbon and inert gases supplied to the vacuum chamber such that the amount of carbon atoms in the layer of C:H which are derived from PECVD of the hydrocarbon gas is less than about 50 at. %.\nIn accordance with embodiments of the present invention, step (c) includes separately controlling the flow rates of each of the hydrocarbon and inert gases to the vacuum chamber such that the amount of carbon atoms in the layer of C:H which are derived from PECVD of the at least one hydrocarbon gas is at least about 30 at. %; step (c) further includes supplying the interior space of the vacuum chamber with at least one hydrocarbon gas of formula CxHy, where x=an integer from 1 to 5 and y=an integer from 2 to 10 and at least one inert gas selected from the group consisting of He, Ne, Ar, Kr, and Xe.\nAccording to particular embodiments of the present invention, step (c) includes supplying the interior space of the vacuum chamber with at least one of acetylene (C2H2) and ethylene (C2H4) as the at least one hydrocarbon gas and Ar as the at least one inert gas; step (a) comprises providing an elongated, cylindrical carbon sputtering target rotatable about its axis of elongation; step (c) further comprises rotating the cylindrical carbon sputtering target about the axis of elongation; and step (b) further comprises applying a bias voltage within the voltage range from 0 to about xe2x88x92150 V to the substrate during step (c).\nEmbodiments of the present invention include comprises providing a magnetic or magneto-optical (MO) recording medium as the substrate in step (b), the surface thereof comprising the exposed surface of an uppermost layer of a stack of layers comprising the medium; and step (c) comprises forming a protective overcoat layer of the hard, abrasion and corrosion-resistant C:H on the exposed surface of the uppermost layer of the medium.\nAccording to particular embodiments of the present invention, step (b) comprises providing a disk-shaped substrate; and step (c) comprises forming the protective overcoat layer to a thickness not greater than about 30 xc3x85.\nStill another aspect of the present invention is a recording medium, comprising:\n(a) a substrate;\n(b) a stack of thin film layers on the substrate; and\n(c) a protective overcoat layer on an uppermost layer of the stack of thin film layers, the protective overcoat layer comprising a novel hard, abrasion and corrosion-resistant material comprised of hydrogenated carbon (C:H) formed by a process comprising simultaneous sputter and plasma-enhanced chemical vapor (PECVD) deposition of the novel hydrogenated carbon (C:H) material, wherein the amount of carbon atoms in the C:H material derived from the PECVD component of the process is less than about 50 at. %.\nAccording to embodiments of the present invention, the amount of carbon atoms in the C:H material derived from the PECVD component of the process is at least about 30 at. %; the position of the Raman G-band of the C:H material of the protective overcoat layer is about 1553 cmxe2x88x921; and the film resistance is as high as about 85 k xcexa9.\nIn accordance with particular embodiments of the present invention, the stack (b) of thin film layers comprises a stack of layers for a magnetic or magneto-optical (MO) recording medium; and the substrate (a) is disk-shaped.\nAccording to specific embodiments of the invention, the stack (b) of thin film layers comprises a stack of layers for a magnetic recording medium; the protective overcoat layer (c) is not greater than about 30 xc3x85 thick; and the stack (b) of thin film layers includes at least one ferromagnetic layer comprising Co.\nA further aspect of the present invention is an apparatus for performing simultaneous sputter and plasma-enhanced chemical vapor (PECVD) deposition of a layer of a novel hard, abrasion and corrosion-resistant, hydrogenated carbon (C:H) material on a surface of a substrate, comprising:\n(a) a vacuum chamber defining an interior space;\n(b) a carbon sputtering target in the interior space of the vacuum chamber, the carbon target being in the form of an elongated cylinder rotatable about its axis of elongation;\n(c) substrate mounting means for mounting a substrate in the interior space of the vacuum chamber, such that a major surface of the substrate is in parallel, facing relation to the elongated cylinder sputtering target; and\n(d) gas supply means for supplying the interior space of the vacuum chamber with at least one hydrocarbon gas and at least one inert sputtering gas at separately controllable flow rates.\nAccording to embodiments of the present invention, substrate mounting means (c) comprises means for mounting a disk-shaped substrate such that a major surface thereof is in parallel, facing relation to the cylindrical sputtering target; and the apparatus further comprises a bias voltage applying means (e) for applying to the substrate a bias voltage within the range from 0 to about xe2x88x92150 V.\nA still further aspect of the present invention is a magnetic recording medium, comprising:\n(a) at least one ferromagnetic thin film layer containing Co; and\n(b) means for protecting the at least one Co-containing ferromagnetic thin film layer from corrosion under high temperature, high humidity environments.\nAdditional advantages and aspects of the present invention will become readily apparent to those skilled in the art from the following detailed description, wherein embodiments of the present invention are shown and described, simply by way of illustration of the best mode contemplated for practicing the present invention. As will be described, the present invention is capable of other and different embodiments, and its several details are susceptible of modification in various obvious respects. Accordingly, the drawings and description are to be regarded as illustrative in nature, and not as limitative."}
-{"text": "X-ray detectors for use in medical computerized tomography are well known in the art. Medical x-ray detectors, such as described and claimed in commonly assigned U.S. Pat. No. 4,031,396 to Whetten et al., are subject to the resolution-limiting effects of k-band x-ray fluorescence. This phenomenon relates to the emission of x-ray photons from the incident x-ray energy-excited gas molecules of the detecting medium and is fully explained in the above-cited patent. The atoms of a gaseous detecting medium, upon interacting with the incident x-ray energy, may themselves emit low energy x-ray photons. When the detecting medium allows a long absorption length for a fluorescent x-ray photon, the photon thus generated by secondary emission may be detected within the detector at a location remote from the point where the incident x-ray entered the detector. Such event degrades spatial resolution and decreases detector efficiency since a portion of the incident x-ray energy is thus not usefully detected. To address this problem, the Whetten et al. patent discloses a detector embodiment which employs voltage and collector plates comprising metals of high atomic number, such as tungsten. Such plates are substantially opaque to x-ray energy and effective to absorb the fluorescent x-rays, thereby minimizing their adverse effect on resolution.\nThe use of tungsten plates, however, causes a decrease in operating efficiency of an x-ray detector. The detector operating efficiency may be expressed as the percentage of incident x-ray energy that is sensed by the detector. It will be appreciated that absorption of fluorescent x-rays by tungsten plates aborts the liberation of charged particles in the detecting medium which can be detected at the collector plates. Hence, that portion of incident x-ray energy responsible for the emission of fluorescent x-rays goes undetected.\nAnother drawback to using tungsten plates is that they can absorb incident x-ray energy just as well as fluorescent x-rays. In the above-cited Whetten et al. patent, the tungsten plates are oriented parallel to the incident x-ray beam. While these plates are quite thin (0.05 mm), x-ray energy striking a plate edgewise will be absorbed. Therefore, in the aggregate, the sum of the thicknesses of the plurality of voltage and collector plates used in the Whetten et al. detector represents a significant incident x-ray absorption area and materially contributes to loss of detector efficiency.\nTypical prior art medical x-ray detectors, such as disclosed in the Whetten et al. patent, employ a relatively low pressure gaseous detecting medium having a correspondingly low density. As a result, the absorption lengths of both the incident x-rays and fluorescent x-rays are relatively long. Such long absorption lengths necessitate the use of tungsten plates which extend a considerable distance into the gaseous detecting medium, e.g. approximatey one inch. The resulting large area presenting by each plate may cause a number of problems. First, microphonics are produced by any motion of the voltage plates with respect to adjacent collector plates. This motion may occur as a result of detector rotation in computerized axial tomography scanning, or even as a result of vibrations transmitted through the building structure supporting the equipment. Such voltage plate movements cause changes in the relative capacitances between the plates, which induce signals in the collector plates. The induced signals are superimposed on the true signals associated with the x-ray image and thus tend to distort spatial resolution. Moreover, the larger plate size must be accommodated, structurally, in order to maintain it rigid with respect to its adjacent collector plates.\nSecond, the large plates in the tungsten plate detector embodiment disclosed in the Whetten et al. patent are arranged in radial fashion aimed at the focal spot of the incident x-ray source. As a result, the position of the detector in relation to the x-ray source must be maintained precisely; otherwise, much of the incident beam will impinge on and be absorbed by the tungsten plates and hence not be detected.\nThird, use of larger voltage plates can result in degradation of kilovolt peak (KVP) linearity. KVP linearity relates in general to a detector's sensitivity to different energy spectra of incident x-rays. The individual collector elements of an x-ray detector are usually calibrated in accordance with a particular x-ray energy spectrum. A patient undergoing examination with a medical computerized tomography x-ray detector may be exposed to an incident x-ray energy spectrum comprising a range of energies. Absorption of x-rays by the patient's body results in a shifting of the average energy of the spectrum toward a higher value. In order to maintain acceptable tomographic image resolution, it is necessary that the individual collector elements respond uniformly to this shifting of the average energy. To this end, it is necessary that the volume of the detecting medium allocated to each collector element be exactly the same. When the respective portions of the detecting medium volume associated with the individual collector elements become unequal, KVP linearity is degraded. Uniformly increasing these volume portions improves KVP linearity. However, in prior art medical detectors, the large tungsten plates require extra care in manufacture to maintain, with precision, equality between the adjacent volume portions respectively associated with the collector elements.\nIndustrial x-ray detectors known in the art are typically employed in such applications as inspection of jet engine turbine buckets, nuclear fuel rods, weldments, and other industrial inspection functions. Typically, these applications require higher resolution and higher energy x-rays than do medical applications. For example, industrial detectors may be used at an average x-ray energy of 300 to 400 KEV while the average x-ray energy in a medical detector application is only on the other of 50 to 125 KEV.\nIn order to achieve the higher resolution required in industrial applications, such industrial detectors have utilized structures in which the individual detector elements are extremely narrow and closely spaced, e.g. individual detector widths of only 2 to 3 mils and spacing of only 1 to 2 mils. Such detector element structure provides a degree of spatial resolution that is neither necessary nor desirable in medical computerized tomography. If such a high resolution detector structure were to be applied in a medical x-ray detector, the lower signal-to-noise ratio that would be experienced would require an intolerable increase in the x-ray dosage to the patient under study in order to achieve acceptable results.\nIndustrial x-ray detectors typically employ a high density gaseous detecting medium in order to minimize the absorption length of the higher energy x-rays being detected. For example, the detector disclosed in U.S. Pat. No. 4,394,578 to Houston et al. employs a detecting medium gas pressure of up to 200 atmospheres.\nThe extremely high (1) detecting medium gas pressures, (2) image resolution structures and (3) x-ray energies are important industrial x-ray detector design considerations which are not applicable to medical x-ray detectors."}
-{"text": "It is not possible to change the arrangement of devices provided in a mobile terminal apparatus in accordance with the direction in which it is held.\nTherefore, there is a risk that the devices may not be used in an appropriate position."}
-{"text": "1. Field of the Invention\nThis invention relates to an electrical energy saving device for motors, which comprises a principal circuit system operating with a control circuit system for automatically adjusting the magnitude of the input power of the motor with respect to the different load states of the motor.\n2. Related Art\nConventionally, all motor control systems have not been provided with a specially designed triggering device for controlling the conducting condition of the motor. Therefore, when the motor is in an idle running state or the load thereof is light, the input current drawn by the motor is still kept at 50-70% of the full-load current. This mode of operation not only consumes too much energy and reduces the power factor of the system, but also decreases the life span of the motor because of the high temperature suffered by the motor when running at this current level.\nAlthough control elements, such as a SCR or a TRIAC, are used in conventional motor control circuits for triggering the running state of the motor, owing to the fact that the magnitude of the starting current is usually 7-8 times of the zero-load current, the specifications of such control elements must be augmented, which causes a significant increase in the manufacturing costs. On the other hand, a timing device is also used in conventional control circuit as a starting and overload protecting device, but it has to be matched with an electromagnetic switch, which increases substantially the volume of the control device as well as the manufacturing costs. Besides, the timing device, which achieves an overload protection only during the starting period, provides no protection during the running state of the motor. Furthermore, conventional motor control circuit systems cannot be installed at the front end of the power source switch, and are very inconvenient in utilization, replacement and adjustment operations."}
-{"text": "Photodiodes operating in the spectral range with a wavelength of up to approximately 1.1 \u03bcm are usually produced from silicon. Photodiodes composed of germanium or InGaAs are sensitive in the infrared range. By using further compound semiconductor materials, for example II-VI semiconductors or else organic semiconductor materials, further spectral ranges can furthermore be covered.\nThe p-n junction in the semiconductor material is produced by using regions having differing doping which provide for a change of charge carrier type. In the region of the p-n junction, a depletion layer, also called space charge zone, depleted of charge carriers is formed on account of diffusion and recombination processes, in which depletion layer an electric field prevails which provides for a separation of the charge carriers in the space charge zone. On account of the continuous thermal generation of charge carriers in the space charge zone and in adjacent regions of the neutral zone from which the charge carriers migrate into the space charge zone as a result of diffusion, a small current, also called reverse current, flows permanently via the space charge zone.\nIf electromagnetic radiation impinges on the space charge zone or, respectively, the adjacent diffusion region, further electron-hole pairs are generated which, in addition to the permanent reverse current, provide for a photocurrent generated by the electromagnetic radiation. In this case, the photocurrent is proportional to the intensity of the electromagnetic radiation. The electromagnetic radiation has to have a minimum energy in order to generate the charge carriers, the known energy being dependent on the band gap of the semiconductor material used. A limit wavelength above which the sensitivity decreases greatly thus results for each semiconductor material. Furthermore, the penetration depth of the electromagnetic radiation is also greatly dependent on the wavelength, the penetration depth increasing with longer wavelengths. In the case of silicon, for example, electron-hole pairs are generated by using visible light up to a limit wavelength of 1.1 \u03bcm, approximately 90% of the incident light radiation being absorbed at a depth of 0.5 \u03bcm in the case of blue light, up to depth of 1.1 \u03bcm in the case of green light and up to a depth of 8.5 \u03bcm in the case of red light.\nIn order to detect the colour information of electromagnetic radiation over a relatively large wavelength range with the aid of photodiodes, as is carried out for example during the detection of colour images with the aid of photodiode arrays in digital cameras, various solutions are known. Thus, for colour resolution, a plurality of photodiodes each coordinated with the spectral range to be detected are interconnected, wherein the individual photodiodes are provided with spectral filters in order to transmit the desired spectral range to the associated photodiode. In this case, the spectral filter can be embodied as an absorption filter for masking out non-relevant radiation and can be used in a Bayer pattern. Furthermore, transmission gratings are known, which are applied over the photodiodes and act as spectral filters in order that only electromagnetic radiation having a specific wavelength is transmitted. Such a construction is described for example by P. B. Catrysser, B. A. Wandell in: Journal of Optical Society of America, Vol. 20, No. 12, December 2003.\nFurthermore, methods are known wherein the electromagnetic radiation is spectrally split with the aid of prisms and fed separately to the assigned photodiodes. Furthermore, there is the possibility, in the case of photodiodes which are sensitive over a wide spectral range, of additionally using discrete filters which are fitted alternately above the photodiode structure in order that the individual spectral ranges are detected sequentially one after another.\nFurthermore, in the case of photodiodes, the space charge zone can be expanded by applying a reverse voltage, in order, with the aid of the larger space charge zone, to improve the effectiveness of the charge carrier detection particularly for long-wave electromagnetic radiation having a large penetration depth. However, the range in which the expansion of the space charge zone can be varied under typical operating voltages for photodiodes is predetermined by the doping concentration and is generally not adjustable over the range of approximately 0.1 to 10 \u03bcm which is of interest for visible light.\nFurthermore, individual photodiodes with spectral resolution are also known wherein no additional spectral filter is required. Thus, U.S. Pat. No. 5,965,875 describes a photodiode construction wherein a plurality of p-n junctions are formed at different depths in the semiconductor material, as a result of which a plurality of space charge zones are formed one above another and separately from one another, which can then be used to simultaneously detect different spectral ranges of the incident light radiation. In this case, each p-n junction has a dedicated electrical connection for reading out the photocurrent detected in the respective spectral range.\nWO 2008/068616 furthermore describes a photodiode construction including a large space charge zone, in which, as a result of transverse electric fields being applied, the charge carriers generated at different absorption depths in the space charge zone are separately collected and read out, such that different spectral ranges of the incident light radiation can be resolved. A similar principle is employed in US 2006/0244089, wherein a large-area electrode is applied above a photoactive layer with a large space charge zone, which electrode is insulated from the photoactive layer and can be used to control the extraction region for the charge carriers in a manner dependent on the absorption depth, in order to obtain a spectral resolution of the incident light radiation.\nThe implementation of spectrally sensitive photodiodes, also referred to as colour sensors hereinafter, as clusters composed of a plurality of photodiodes specifically tailored to the spectrum that is respectively to be detected, or as broadband photodiode with additional colour filters is complex in production and can be integrated only with difficulty into standard processes for producing semiconductor circuits, in particular in CMOS circuits. Furthermore, a large area requirement is generally necessary, which limits the image resolution when such colour sensors are used in CMOS cameras, for example.\nMoreover, generally only a discrete colour resolution is possible with the known colour sensors. This also applies, in particular, to photodiodes having multiple p-n junctions in order to provide different absorption depths and hence spectral sensitivities. Furthermore, each p-n junction has to be individually electrically connected in order that the charge carriers generated can be read out separately and the individual colours can thus be resolved.\nThe design of colour sensors as lateral photodiodes having a multiplicity of electrical connections in order, by using corresponding transverse fields, to read out the charge carriers from different absorption depths and thus to obtain a colour resolution demands a large area requirement and is furthermore complex in production. In the case of the photodiode structure proposed in US 2006/0244089, wherein control is obtained with the aid of an insulated surface electrode, there is the difficulty that the electrode material already absorbs electromagnetic radiation prior to penetration into the actual photoactive layer. As a result, particularly in the short-wave spectral range, less light intensity is available for electrical signal formation.\nFor these and other reasons there is a need for the present invention."}
-{"text": "1. Field of the Invention\nThe present invention relates to an LPI fuel supply system. More particularly, the present invention relates to an LPI fuel supply system that securely supplies an engine with fuel from a fuel tank.\n2. Description of Related Art\nGenerally, a liquefied petroleum injection (LPI) engine uses an injector to inject liquefied fuel into a cylinder.\nThe operation principle of the LPI engine is schematically explained as follows: an electronic control unit (ECU) receives input signals from sensors to detect the condition of the engine and controls a fuel pump, the injector, and an ignition coil so as to optimize air/fuel ratio and engine performance.\nThat is, the fuel pump is controlled according to the fuel amount that is necessary to supply liquefied fuel to the engine, and the injector sequentially injects the fuel into the cylinders to realize an optimized air/fuel ratio.\nThe LPI system that is provided to the LPI engine includes a fuel pump that is disposed inside or outside the fuel tank so as to supply the liquefied fuel to the engine.\nIf the fuel pump supplies the fuel with uniform flux, the injector injects the liquid fuel. Here, the pressure regulator is disposed so as to uniformly sustain the pressure that is applied to the injector. If the pressure that is formed in the supply line is higher than a control pressure, the passage of the pressure regulator is opened to return the fuel to the fuel tank.\nHowever, when the fuel pump is deteriorated or foreign substances are accumulated in the main filter of the fuel pump, there is a problem that the supply capacity of the pump becomes weakened.\nThe information disclosed in this Background of the Invention section is only for enhancement of understanding of the general background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art already known to a person skilled in the art."}
-{"text": "Presently, germination of plant seeds of flowering plants and/or of plant spores from non-flowering plants, may generally involve wetting the seed and/or spore to trigger germination. Such a process may generally involve a given seed and/or a given spore, a moisture source, and a growth medium. Often the moisture source may be water that may have wetted the growth medium, which may then retain some of the water. The growth medium may comprise: soil, potting soil, gardening soil, peat, moss, gravel, sand, vermiculite, fertilizer, manure, paper products (e.g. paper towels, napkins, tissues, and the like), fabrics, textiles, linens, and the like. Once the seed and/or the spore may be brought into contact with sufficient moisture, germination may initiate and then subsequent growth may be dictated by temperature, moisture content of the growth medium, and presence of light/dark cycles, the severity of such light, and possibility an applicable growing zone if not using artificial light and/or temperature sources.\nHowever, the process of initiating seed and/or spore germination has historically and presently remains relatively benign, boring and uninteresting. Aside from wetting the seed and/or the spore in a growth medium, there is little to see and/or do.\nThere then is need in the art for devices, apparatuses, methods, and/or systems of initiating germination in a way which may be exciting, interesting\u2014providing a dynamic visual event to observe and possibly specific smells to smell. For example, there is a need in the art for a solid container with a specific shape, that may comprise one or more seeds (spores), wherein the solid container with the specific shape may be triggered to dissolve in an aqueous acid base chemical reaction producing bubbles and fizzing to expose the plant seeds (and/or plant spores); wherein the wetted plant seeds (and/or wetted plant spores) may be wetted to initiate germination. And in some embodiments, the bubbles and fizzing may also emit specific smells (odors and/or fragrances). The specific shape may be a shape not traditionally associated with growing plants and thus it housing seeds (spores) which may germinate after the specific shape dissolves, may offer a pleasant surprise to the user.\nIt is to these ends that the present invention has been developed."}
-{"text": "High molecular weight polyesters are commonly produced from low molecular weight polyesters of the same composition by solid state polymerization. The low molecular weight polyesters which are used in such solid state polymerizations can be prepared by conventional melt polymerizations. Solid state polymerization is generally considered advantageous in that the handling of high molecular weight ultra-high viscosity molten polymers is eliminated during the solid state polymerization phase. Thermal degradation during the solid state portion of the polymerization is also essentially avoided.\nIn melt polymerizations, the major part of the reaction is transesterification due to the fact that thermal pyrolysis nullifies much of the esterification of carboxyl end groups. A much larger percentage of the reaction in a solid state polymerization is esterification. Thus, a high molecular weight can be obtained in a solid state polymerization with a lower carboxyl concentration than can be achieved in a melt polymerization.\nThe low molecular weight polyester prepolymers utilized in solid state polymerizations are generally in the form of pellets, chips, or finely divided powder. Such pellets can vary greatly in size: however, as a general rule, the smaller the size of the pellets of polyester prepolymer the faster the solid state polymerization with proceed. Such polyester prepolymers are generally converted from the amorphous to the crystalline state prior to solid state polymerization in order to raise their sticking temperature. This is done in order to keep the pellets or chips of polyester prepolymer from sticking together as a solid mass in the solid state polymerization reactor.\nIn the solid state polymerization of a polyester prepolymer the polymerization is carried out at an elevated temperature which is below the melting point of the polyester resin. Such polymerizations are normally conducted in the presence of a stream of inert gas or under a vacuum. Solid state polymerizations are normally conducted on a commercial basis in the presence of a stream of inert gas since it serves to remove volatile reaction products and helps to heat the polyester.\nHeretofore, the form of the polyester prepolymer has essentially dictated the type of solid state polymerization process which could be employed in order to convert the low molecular weight polyester prepolymer into high molecular weight polyester resin. For example, it has generally been accepted practice to use polyester prepolymer in the form of pellets or chips in vacuum and static bed processes and finely ground powder in fluidized bed processes. The reason for this is that experience has shown that finely ground powders tend to agglomerate in vacuum processes, resulting in slower polymerization rates and a need to regrind the high molecular weight polyester resin produced. Experience has also shown that, in static bed processes finely ground powders will channel or fissure, resulting in uneven polymerization and prolonged polymerization rate. On the other hand, the use of pellets or chips in fluidized bed processes is not economically feasible in view of the velocity and volume of inert gas needed to suspend the pellets or chips and the size of the equipment required to do so.\nPolyester prepolymers which are in the form of finely divided powders solid state polymerize at faster rates than do polyester prepolymers which are in the form of pellets or chips. However, polyester prepolymers which are in powder form are difficult to handle and generally must be polymerized in fluidized bed processes. Additionally, the high molecular weight polyester resins which are made utilizing prepolymers which are in powder form are also in the form of powders which are more difficult to process into articles of manufacture. For these reasons polyester prepolymers in powder form have not been widely utilized in commercial solid state polymerization techniques."}
-{"text": "The means with which user of computing devices can communicate with each other continues to grow. Where individuals may once have been reached at one or two telephone numbers, users may presently be reached via, for example, multiple telephone numbers (e.g., home, work, and mobile), home and work fax numbers, multiple email addresses, personal and professional web pages, and various social networking profiles (e.g., Facebook, Twitter, etc.).\nGiven the amount of contact information individuals currently maintain, a need exists for a new approach to managing, organizing, and presenting the various ways individuals may be contacted."}
-{"text": "Media networks may simultaneously broadcast multiple types of programming streams, such as high definition (HD) and standard definition (SD) programming streams. In general, HD programming may have a 16:9 aspect ratio and a resolution of about 1920 vertical lines of resolution by 1080 horizontal lines of resolution, or 1920\u00d71080 (2.1 megapixels per frame) or 1280\u00d7720 (0.9 megapixels per frame). Standard definition (SD) is generally used to refer to signal formats that are lower than HD. Media networks may provide various types of content in the available formats based on user demand, content provider request, network requirements. The different formats require different resources, for example, HD content requires more bandwidth than SD bandwidth. In general, media networks increase resource utilization in order to provide higher profile content, such as transmitting HD content using additional bandwidth. However, increasing resource utilization to accommodate content introduces inefficiencies and decreases the profitability of providing multiple types of programming formats. Accordingly, media networks would benefit from methods configured to efficiently allocate bandwidth for content for transmission to viewers using existing resources."}
-{"text": "A conventional position detection method typically involves measuring a position of an object using signals transmitted from a satellite, such as a GPS.\nAnother known position detection method involves receiving signals transmitted from a plurality of access points at a terminal, and calculating the position of the terminal from differences in times of reception of the signals.\nMore specifically, a method for detecting the position of a mobile terminal in a cellular phone system has been proposed which includes calculating the differences in the reception times of signals transmitted from the access points to the mobile terminal (e.g. differences in propagation times of the signals from the respective access points to the mobile terminal T1\u2212T2 and T3\u2212T2), and multiplying the propagation time differences by light velocity to calculate differences in propagation distances of the signals from the mobile terminal to the respective access points, thereby determining the position of the terminal, in the following manner (see patent document 1).D1\u2212D2=c(T1\u2212T2)D3\u2212D2=c(T3\u2212T2)\nAnother method for detecting a position of a node in a wireless LAN system has been proposed which includes calculating differences in reception times of a signal (differences between the reception times at the respective access points Ti\u2212Tl), which has been transmitted from the node (terminal) and received at the access points, and multiplying the reception time differences by light velocity to calculate differences between propagation distances of the signal from the node to the respective access points ({[P\u2212Pi]\u2212[P\u2212Pl]}=c (Ti\u2212Tl), i=2, . . . , n), thereby detecting or determining the position of the node, in the following manner (see patent document 2).    [Patent document 1] JP-A No. 181242/1995    [Patent document 2] JP-A No. 101254/2004"}
-{"text": "Many people desire to have easy access to their wireless devices at any time\u2014even when exposed to cold environments. Accordingly, there is a need for improved products for allowing users easy access to their wireless devices."}
-{"text": "1). Field of the Invention\nThe field of this invention is directed to whistles designed to be placed within the mouth of a human and used to produce a high pitched shrill sound by blowing air through the whistle.\n2). Description of the prior Art\nThe construction of whistles has long been known also the construction of whistles out of a confectionery product, such as hard candy, has long been known. However, heretofore, whistles constructed of the prior art have not been designed in such a way as to maximize the creation of a high pitched shrill sound.\nWhistles are commonly used as a self entertainment device. However, whistles do have utility in that they may be used as a means of getting someone's attention located at some distance away from the whistler. Also, whistles can be used as a means of instructing animals to do certain things, such as calling cows to a barn.\nA primary objective of the present invention is to construct a whistle of superior quality than previously constructed whistles that are able to produce an extremely loud high pitched shrill sound.\nAnother objective of the present invention is to construct a whistle which can be readily manufactured from either a permanent type of material, such as plastic, or can be manufactured from an edible material such as a hard candy.\nAnother objective of the present invention is to construct a whistle which can be manufactured inexpensively and thereby sold to the ultimate consumer at an inexpensive price.\nThe whistle of the present invention is constructed of a body which has a top wall and a bottom wall which are interconnected by a sidewall. An air entry through hole is formed between the top wall and the bottom wall with this air entry through hole being located directly adjacent the sidewall. Formed within the body is an air exhaust chamber which is tapered and has straight walls connected between a narrow tip and a widened base. The air entry through hole connects with the narrowed tip. The widened base forms an air exhaust opening within the top wall. The air exhaust chamber is slantingly located within the body permitting a support stick to be mounted within the body. The portion of the body including the air entry hole is to be located within a mouth of a user, and by forcing air into the air entry opening and into the air exhaust chamber a high pitched shrill sound is produced."}
-{"text": "1. Field of the Invention\nThe present invention relates to an optical component such as a polarizing beam splitter or a phase plate in which a multilayer film is provided on a surface of a substrate to polarize incident light, and in particular, to an optical component having a multilayer film with a high aspect ratio as well as a method of manufacturing the same.\n2. Related Background Art\nHitherto, an optical component such as a lens, used in an electronic equipment is based on an optical interference and has thin layers stacked on a surface of an optical element as an antireflection film. Further, on the basis of a similar optical interference, thin layers are stacked on a surface of a substrate composed of a transparent optical element to polarize an incident light, thereby forming a polarizing beam splitter or a phase plate. Such an interference film is inexpensive and has good reproducibility, and is thus put to practical use widely.\nFIG. 5 shows an example of a conventionally used interference film cubic type polarizing beam splitter. In FIG. 5, reference numerals 12a and 12b denote transparent body substrates each having a surface inclined at 45\u00b0. Reference numerals 13 and 14 denote a low refractive index layer and a high refractive index layer, respectively. Further, the arrows shown in the figure indicate the direction in which a light travels. A normal light has a S polarization beam that is a vibration component having a light vibrating direction perpendicular to the drawing of FIG. 5 and a P polarization beam that is a vibration component having a light vibrating direction parallel to the drawing of FIG. 5. The polarizing beam splitter can split an incident light into two components as shown in the figure by allowing the P polarization beam to pass therethrough while reflecting the S polarization beam.\nHowever, although the cubic type polarizing beam splitter shown in FIG. 5 is very easy to manufacture and has good reproducibility, there is a problem as to its optical characteristics that the reflection and transmission characteristics change markedly depending on the incident angle of light.\nFor example, when the size of a device such as a liquid crystal projector is to be reduced, its optical path length must be shortened. The use of the conventionally used interference film cubic type polarizing beam splitter enables the optical path length to be shortened. However, because the angle at which a light is incident on the polarizing beam splitter varies markedly depending on the position of the incidence, the reflection or transmission characteristics are not constant, thus making it impossible to obtain uniform spectral characteristics within a plane. Therefore, to further improve the functions of electronic equipments, there is a need for an optical component such as a polarizing beam splitter or a phase plate that provides good optical characteristics, regardless of the incident angle of light.\nTo meet this need, a polarizer has recently been proposed which polarizes an incident light utilizing the anisotropy of propagation characteristics of the periodic structure of a transparent body and in which high refractive index layers and low refractive index layers are repeatedly stacked alternately on each other to provide a groove form that is bent periodically at a pitch less than the wavelength of visible light as disclosed in Japanese Patent Application Laid-Open No. 2000-56133. This publication describes the use of Si, GaAs, TiO2, and TaO2 as the high refractive index layer and the use of SiO2 as the low refractive index layer. Further, this publication describes an example in which the high refractive index and the low refractive index layers has a periodic thickness of 0.32 \u03bcm and in which the groove has a pitch of 0.4 \u03bcm and a depth of 0.2 \u03bcm.\nWith the configuration described in Japanese Patent Application Laid-Open No. 2000-56133 above, the polarizer can effectively function as a polarizing beam splitter for a light of a wavelength of 1 \u03bcm. However, the polarizer is difficult to use if an incident light to be polarized is a visible light having a wavelength of 400 to 700 nm. That is, Si and GaAs that constitute the high refractive index layer has a low transmissivity for visible light. Thus, the polarizer is difficult to use in an optical component such as a liquid crystal projector which must allow a light of a visible range to pass therethrough. Further, because the difference between the refractive index of TiO2 and TaO2 constituting the high refractive index layer and that of SiO2 constituting the low refractive index layer is too small, sufficient polarization characteristics cannot be obtained.\nOPTICS LETTERS (Rong-Chung et al., Vol. 21, No. 10, p. 761, 1996) describes a polarizing beam splitter comprising a plurality of dielectric multilayer films arranged in lines on a transparent substrate as shown in FIG. 6. In FIG. 6, reference numeral 21 denotes an optical component such as a beam splitter or a polarizing plate. Reference numeral 22 denotes a transparent substrate as an optical element, and reference numerals 23 and 24 denote low refractive index and high refractive index layers, respectively. The low refractive index layers 23 and the high refractive index layers 24 are alternately stacked on a surface of the transparent substrate 22 to form an HL alternate layer 25. Reference numerals 26a, 26b, 26c, 26d denote dielectric multilayer films each formed in a line. It is known that the line-shaped dielectric multilayer films 26a, 26b, 26c, 26d which form periodic recesses and protrusions function as a diffraction grating and have excellent characteristics as an optical component that polarizes an incident light. Specifically, even if an incident light is a visible light having a wavelength of 400 to 700 nm, good polarization characteristics can in principle be obtained by making the line width of each of the line-shaped dielectric multilayer films less than the wavelength of the visible light, preferably 0.1 \u03bcm or less. Further, for its manufacturing method, a patterning technology for a super LSI can be applied to manufacture the periodic recesses and protrusions using dielectric multilayer films. Examples of known patterning technology include wet etching using an etchant such as hydrofluoric acid, dry etching such as reactive ion etching, or the like.\nHowever, when a polarizing beam splitter is manufactured using the line-shaped dielectric multilayer films shown in FIG. 6, a wet process will be performed in an etching step or washing step, as described above. Further, a drying step is required after this wet process. Thus, a processing liquid flowing into the recesses during the above steps pulls the protrusions by a capillary force, so that a stress is applied to the dielectric multiplayer films as the protrusions to incline them. This tendency is significant when in the recesses and protrusions formed by the line-shaped dielectric multilayer films, the thickness of the dielectric multilayer film is relatively greater as compared with the width of the recess, i.e., the aspect ratio of the recesses and protrusions is high. This tendency becomes very significant particularly when the aspect ratio is 4 or more.\nAs a manufacturing method for avoiding this, Japanese Patent Application Laid-Open No. 2001-165568 discloses a drying method utilizing a supercritical state. According to this method, a processed pattern is immersed into normal hexane as a nonpolar solvent, and then sealed together with liquid carbon dioxide in a reaction chamber, so that the normal hexane is replaced with the liquid carbon dioxide to dry the processed pattern. This method is effective in preventing defects in the pattern. However, the inner pressure of the reaction chamber must be set to 7.5 MPa, thus requiring a high pressure container. Further, an organic solvent is required as a nonpolar solvent, which may pollute the environment."}
-{"text": "This invention is in the field of circuitry for driving light-emitting diodes (LEDs). Embodiments are more specifically directed to LED driver circuitry in transmitters of photoplethysmography (PPG) systems.\nPhotoplethysmography (PPG) is a modern technology that has proven useful for the measurement of cardiovascular function in humans. According to this technology, fixed wavelength light from a light-emitting diode (LED) is emitted into the skin of a human subject, and is sensed by a photodiode (PD) after transmission through the skin and underlying tissue. The characteristics of the sensed light allows measurement of medical parameters such as oxygenation, pulse rate, respiratory function, and the like.\nConventional PPG sensors include the well-known pulse oximeter, for example of the type that clip-on onto the finger of the patient. Pulse oximeters typically measure the oxygen saturation of circulating blood from a comparison of the absorption of light in the dermis and subcutaneous tissue at two different wavelengths. So-called \u201cwearable\u201d devices such as heart rate sensors also utilize this technology, but need only measure light absorption at a single wavelength.\nFIG. 1 generically illustrates the architecture of a conventional PPG system. Transmitter 3 in this system includes LED 2, which has its anode biased by the Vdd power supply voltage and its cathode coupled to ground via LED driver 4. As mentioned above, PPG pulse oximeters will included multiple LEDs 2 (e.g., red and infrared), operated in time-multiplexed fashion. When forward-biased by LED driver 4, LED 2 emits light into the patient, for example the index finger of the patient. Receiver 7 includes photodiode 6, which has its cathode biased at the Vdd power supply voltage and its anode connected to the input of amplifier 8, and which is normally reverse-biased so that photons impinging photodiode 6 will produce a current detectable by amplifier 8. In this manner, photodiode 6 senses the extent to which the light emitted by LED 2 is transmitted through the subject. The output of amplifier 8 is forwarded to the desired processing and analysis circuitry of the PPG system to determine the desired medical measurement, such as the oxygenation of the patient's blood.\nThe ability of any PPG system to accurately and precisely measure the parameter of interest is based on the signal-to-noise ratio (SNR) of the overall system, considering both its transmitter and receiver. For example, it has been observed that an SNR of at least about 30 dB, for the PPG system as a whole including the transmission channel through patient tissue, is necessary in order to measure pulse rate to an accuracy of 1 beat per minute (bpm). A complicating factor in practice is that the system SNR depends on the perfusion index of the patient, as illustrated in FIG. 2. Perfusion index is the ratio of the AC signal due to pulsatile blood flow to the DC background level of the light signal passing through the patient's peripheral tissue, and depends largely on the health and physical condition of the patient. As such, for the PPG system of FIG. 1, the perfusion index is reflected in the amplitude of the AC pulses of the received light (e.g., as output by amplifier 8) relative to its DC level. As shown in FIG. 2, the system SNR increases at higher perfusion index values.\nConversely, if the SNR of the PPG system can be increased, the system can measure the pulse rate and blood oxygenation in a wider range of patients, particularly those of poorer health and thus lower perfusion indices. FIG. 2 shows two SNR vs. perfusion index plots 9a, 9b. Plot 9a illustrates the relationship of system (transmit\u2014channel\u2014receive) SNR to perfusion index for the case in which the transmitter SNR is 95 dB; as known in the art, to attain the required system SNR, the SNRs for the transmitter and receiver must both be higher than that required system SNR. At that transmit SNR, the PPG system is able to measure pulse rate to an accuracy of one bpm only for patients exhibiting a perfusion index of at least about 0.06. In contrast, plot 9b illustrates that if the transmitter is able to operate at an SNR of 110 dB, pulse rate measurements at an error of 1 bpm can be made for patients with perfusion index values as low as about 0.01. Accordingly, noise in the transmitter of the PPG system is a critical factor in covering a wide range of patients.\nSo-called \u201cwearable\u201d electronic devices, such as fitness monitoring devices, have recently become popular. In addition to fitness monitoring devices, wearable medical monitoring devices are being contemplated for use in healthcare, for example to monitor the recovery or progress of a patient suffering from a medical condition. As such, the use of PPG to obtain oxygenation, pulse rate, and other measurements by way of a wearable device, particularly such a device that can be worn all day, is desirable. In this context, battery life becomes of critical performance.\nIn this regard, an important electrical parameter of a transmitter in a battery-powered system, such as a wearable device, is the \u201cheadroom\u201d of the LED driver. As well-known by those in the art, it is desirable that battery-powered systems operate at low power supply voltages to reduce power consumption and to reduce the cost of the battery itself. In transmitter 3 of FIG. 1, the voltage drop across LED 2 in its operating state is defined by its material. The headroom, shown as Vhead in FIG. 1, is the voltage required by LED driver 4 beyond the LED voltage drop. Conversely, the minimum Vdd power supply voltage is the sum of the voltage drop across LED 2 and the headroom Vhead of LED driver 4. For battery powered systems, therefore, it is desirable that the headroom Vhead required by LED driver 4 be minimized, especially considering that the output voltage from conventional batteries tends to sag over time.\nAs mentioned above, it is desirable to minimize power consumption in PPG systems, particularly those in battery-powered wearable devices intended for \u201call-day\u201d use. Duty cycling of LED driver 4 in transmitter 3 is a common approach to reducing system power consumption. It is therefore desirable for LED driver 4 to exhibit fast switching, and rapid settling times, so that the \u201con\u201d pulse width can be reduced as much as possible and thus minimizing power consumption.\nFIGS. 3a through 3c illustrate examples of conventional LED drivers for PPG systems. The circuit of FIG. 3a is an example of a typical LED driver circuit, such as used in a PPG system as described above. In this circuit, LED 10 has its anode at the Vdd power supply voltage and its cathode connected to the drain of rise time control n-channel MOS transistor 12. Transistor 12 has its source connected to the drain of n-channel driver transistor 14, which is connected to ground via variable resistor 16. Resistor 16 operates to control the current drawn through LED 10 and transistor 12, at a resistance typically set by a digital-to-analog converter (DAC). The gate of transistor 12 receives a control voltage from rise time controller 13, which is an adjustable circuit block that controls the conduction of transistor 12 to attain the desired rise and fall times in the turn-on and turn-off characteristics of LED 10. Amplifier 18 receives a reference voltage VREF at its non-inverting (positive) input and a feedback voltage from the source of transistor 14 at its inverting (negative) input. Output voltage VGATE from amplifier 18 is applied to the gate of driver transistor 14. According to this arrangement, amplifier 18 operates to drive the gate voltage VGATE at driver transistor 14 so that reference voltage VREF at the source node of transistor 14. Reference voltage VREF is modulated to selectively forward bias LED 10.\nThe LED driver circuit of FIG. 3a provides certain advantages in a PPG system. Specifically, variable resistor 16 tends to reduce the transmitter noise in this circuit, and the ripple exhibited by this circuit is also quite low. However, it has been observed that this arrangement is vulnerable to significant input flicker noise from amplifier 16, degrading transmitter performance. The settling time of this LED driver is also quite slow, due to the bandwidth of amplifier 18. In addition, the LED driver of FIG. 3a is not conducive to implementation in low voltage, battery-powered, applications because of its large headroom voltage, specifically the sum of the drain-to-source overdrive voltages of transistors 12 and 14 plus the voltage drop across resistor 16.\nFIG. 3b illustrates a conventional LED driver with very low headroom requirements as useful for a PPG system. In this circuit, power supply 20 applies the Vdd bias to LED 10 through inductor 22. N-channel driver transistor 26 has its source-drain path connected in parallel with LED 10 between inductor 22 and ground. The gate of transistor 26 receives the output of pulse-width modulator (PWM) 24. During the \u201coff\u201d pulses, transistor 26 shunts the inductor current through inductor 22 to ground; during the \u201con\u201d pulses, the Vdd power supply voltage forward biases LED 10, such that the inductor current is conducted through LED 10 to ground. FIG. 3b illustrates the behavior of output current Iout through LED 10 over a sequence of pulses from PWM 24, illustrating that the output current Iout appears as a sequence of triangle waves. While the headroom required by this LED driver is quite low, it has been observed that transmitter noise is quite high in this arrangement, which reduces the patient coverage as discussed above relative to FIG. 2. In addition, significant ripple is present in the LED driver of FIG. 3b. \nIn the circuit of FIG. 3c, LED 10 has its anode at the Vdd power supply voltage and its cathode connected to the drain of n-channel MOS transistor 14; the source of transistor 14 is at ground. A reference voltage VREF is applied to the gate of transistor 14, and is modulated to turn LED 10 on and off, thus controlling the emission of light. This simple driver of FIG. 3c has a low headroom voltage of only the drain-to-source voltage overdrive of transistor 14, and exhibits no ripple. However, because the driver of FIG. 3c is quite noisy, its use in the transmitter of a PPG system will have limited patient coverage, as discussed above relative to FIG. 2.\nBy way of further background, auto-zeroing techniques for removing offset voltage and drift of operational amplifiers are known in the art, as described in Kugelstadt, \u201cAuto-zero amplifiers ease the design of high-precision circuits\u201d, Analog Applications Journal, 2Q 2005 (Texas Instruments Incorporated), pp. 19-28, incorporated herein by reference."}
-{"text": "1. Field of the Invention\nThe present invention relates to small, chain-driven carts (sometimes referred to as \"go-carts\"), and the like having an internal combustion engine of, for example, five-ten (5-10) horse power that drives the cart rear wheels via an axle and centrifugal clutch. Even more particularly, the present invention relates to an improved engine shaft, clutch, and chain-driven axle assembly for powering small motorized carts wherein an improved lubrication arrangement allows grease to be injected through an engine drive shaft (or a short intermediate shaft or \"jack\" shaft) of the cart via a porting system so that the clutch is lubricated when rotating at high speeds, namely, when the cart is at rest and the engine idles.\n2. General Background\nSmall chain-driven motorized carts are commercially available from a number of manufacturers. These carts are used primarily for off road motoring at relatively low speeds. They are frequently used for backyard and off road family fun, and can be found at amusement parks and fairs on occasion.\nThese small chain-driven carts include a metal frame that is generally of welded steel, using square and round tubing and plate steel, for example. Such small carts use an internal combustion engine for power, usually on the order of five to ten (5-10) horse power, as an example. The cart frame is constructed very low to the ground for stability purposes. The engine on such a cart is usually a commercially available internal combustion engine, having a drive shaft that carries a centrifugal clutch. A sprocket on the clutch engages an endless belt or chain. Another sprocket rigidly affixed to the axle also engages the endless belt or chain. When the engine of the cart idles, the centrifugal clutch by design is disengaged. When the user presses a throttle (e.g., foot pedal) the engine revolutions increase and the clutch engages so that the sprocket of the clutch rotates and drives the endless belt or chain and its associated axle sprocket and attached axle.\nBy placing the cart low to the ground, equipping the cart with relatively small wheels, and by using the cart in a off road environment, the centrifugal clutch is subjected to dirt, sand, dust, and debris which accumulates in and around the clutch assembly. This produces an abrasive environment which can quickly wear out the clutch parts. The centrifugal clutch is typically rotated at high speeds when the cart idles. The clutch engages upon increasing engine revolutions so that centrifugal members engage the clutch and it forms a rigid connection with the drive shaft of the cart.\nTypically, these carts use a horizontal shaft engine. The drive shaft extends horizontally from one side of the engine. However, it should be understood that a short separate intermediate drive shaft sometimes referred to a \"jack-shaft\" can be used to drive the carts. A \"jack-shaft\" is simply a short shaft mounted between bearings and in spaced-apart position away from the horizontal shaft of the engine. The \"jack-shaft\" is driven by the motor shaft using a belt or sprockets and a chain, for example. In either situation, the present invention is designed to provide an improved lubrication arrangement for the engine shaft, and/or a \"jack-shaft\" used with the engine drive."}
-{"text": "1. Field of the Invention\nThe invention pertains to the field of hearing protection. More particularly, the invention pertains to hearing protectors having talk-through capability.\n2. Description of Related Art\nTalk-through listening systems typically comprise hearing protection devices (i.e. earplugs, headsets, helmets) with microphones mounted on their outside surface. The microphones convert sound pressure levels (SPLs) to an electrical signal. The electrical signals are amplified and used to drive audio transducers such as speakers in headphones, helmets, or earplugs or the like. Electronic automatic gain control and limiting devices ensure that the SPL generated in the canals of the user are kept within safe levels. In this way, the user can monitor sound in a high-SPL environment at safe SPLs.\nOne example of such an environment is the flight deck of aircraft carriers. Aircraft maintainers need to service aircraft under extremely high noise conditions. However, they need to be aware of their surroundings and must communicate to each other. The maintainers must know if a vehicle or plane is approaching from areas outside of their view. For this, they use acoustic cues. Moreover, the maintainers often need to communicate with each other face-to-face in high acoustic noise, without the use of radios."}
-{"text": "This invention relates generally to surgical wound retractors and more specifically to circular wound enactors that provide circumferential retraction of an incision or wound and also provide isolation of the surgical incision or wound from the lumen of the retractor.\nRetraction and isolation of a surgical incision can be an important element associated with a surgical procedure. Adequate access may be provided by a circumferential retraction device so that a surgeon may operate in a clear and open, generally unrestricted field. In addition, a circumferential retractor having a fluid impermeable sleeve extending through the incision or wound may prevent excessive bleeding or contamination of adjacent tissue. In some surgeries, it is important to avoid contaminating the incision or wound site.\nCommon procedures for placing current circumferential wound retraction devices involve placement of an internal retention ring that is coupled to a thin-walled sleeve or tube that is further coupled to an external retention member. The external retention member is wound or turned upon itself to wind the sleeve upon the external retention member. The physics associated with winding a sleeve upon a circular winding collar or external retention member dictate that ample energy must be applied to the external winding motion to provide adequate circumferential retraction of a linear incision. The initial windings of the process are not overwhelming, however, as the process continues the retraction becomes more difficult and the mechanical advantages first appreciated begin to dissipate. The final windings are more difficult since a balance between the force required to wind the sleeve upon the external retention member must be balanced with the force acting to unwind the sleeve from the external retention member.\nThere remains a need to provide a substantially rigid, noncompliant circumferential retractor that easily shortens the length of an associated tubular sleeve coupled to an internal retention member so that the shortening of the sleeve results in dramatic retraction of an incision or wound. Additionally, the retention members must remain substantially circular so that the incision is retracted symmetrically."}
-{"text": "In wireless communication systems, one base station (BS) generally provides services to a plurality of user equipments (UEs). The BS schedules user data for the plurality of UEs, and transmits the user data together with control information containing scheduling information for the user data. In general, a channel for carrying the control information is referred to as a control channel, and a channel for carrying the user data is referred to as a data channel. The UE finds control information of the UE by searching for the control channel, and processes data of the UE by using the control information.\nIn order for the UE to receive user data assigned to the UE, control information for the user data on a control channel must be received. In a given bandwidth, a plurality of pieces of control information for a plurality of UEs are generally multiplexed within one transmission interval. That is, to provide a service to the plurality of UEs, the BS multiplexes the plurality of pieces of control information for the plurality of UEs and then transmits the control information through a plurality of control channels. The UE searches for control channel of the UE among the plurality of control channels.\nBlind decoding is one of schemes for detecting specific control information from the plurality of pieces of multiplexed control information. The blind decoding attempts to recover a control channel by using several combinations of information in a state where a UE has no information required to recover the control channel. That is, in a state where the UE does not know whether control information transmitted from the BS is control information of the UE and the UE does not know in which portion the control information of the UE exists, the UE decodes all pieces of given control information until the control information of the UE is found. The UE can use information unique to each UE to detect the control information of the UE. For example, when the BS multiplexes control information of each UE, an identifier unique to each UE can be transmitted by being masked onto a cyclic redundancy check (CRC). The CRC is a code used for error detection. The UE de-masks unique identifier of the UE from the CRC of the received control information, and then can detect the control information of the UE by performing CRC checking.\nIf the UE cannot correctly detect the control information of the UE from the pieces of multiplexed control information, user data on the data channel cannot be decoded. Therefore, fast and correct detection of the control information has a significant effect on overall system performance. However, it may be difficult to detect the control information when using only blind decoding. Since each UE may require different control information and may use a channel encoding scheme by the use of a different code rate, each UE may have control information with a different size. Therefore, the number of blind decoding attempts may be significantly increased in a control region in which control information is transmitted. Battery consumption of the UE increases in proportion to the number of detection attempts.\nMeanwhile, as a mobile communication system of a next generation (i.e., post-3rd generation), an international mobile telecommunication-advanced (IMT-A) system is standardized aiming at support of an Internet protocol (IP)-based seamless multimedia service in an international telecommunication union (ITU) by providing a high-speed data rate of 1 gigabits per second (Gbps) in downlink communication and 500 megabits per second (Mbps) in uplink communication. In a 3rd generation partnership project (3GPP), a 3GPP long term evolution-advanced (LTE-A) system is considered as a candidate technique for the IMT-A system. The LTE-A system is evolved to increase a completion level of the LTE system, and is expected to maintain backward compatibility with the LTE system. This is because the provisioning of compatibility between the LTE-A system and the LTE system is advantageous in terms of user convenience, and is also advantageous for a service provider since existing equipment can be reused.\nAccordingly, there is a need for a method in which a BS effectively transmits control information and a UE effectively receives data by using the control information in a next generation mobile communication system such as an LTE-A."}
-{"text": "1. Field\nAspects of the present invention relate to a rechargeable battery. More particularly, aspects of the present invention relate to a rechargeable battery having a short circuit member.\n2. Description of the Related Art\nUnlike a primary battery, a rechargeable battery is designed to repeatedly perform charge and discharge functions. Small capacity rechargeable batteries may be used in small portable electronic devices such as mobile phones, laptop computers, and camcorders, and large capacity rechargeable batteries may be used in power sources for driving motors of hybrid vehicles and the like.\nIn recent years, a high power rechargeable battery using a non-aqueous electrolyte having high energy density has been developed. The high power rechargeable battery is a large capacity rechargeable battery with a plurality of rechargeable batteries connected to each other in series to be used to drive a motor of a device requiring high power, such as an electric vehicle.\nThe large capacity rechargeable battery generally includes a plurality of rechargeable batteries which are connected to each other in series, and the rechargeable battery may have a cylindrical shape or an angular shape.\nSuch a rechargeable battery requires a safety apparatus capable of reducing or preventing excessive heat and ensuring safety upon overcharge.\nThe above information disclosed in this Background section is only for enhancing understanding of the background of the invention and therefore it may contain information that does not form prior art."}
-{"text": "1. Field of the Invention\nThis invention relates to a faucet assembly, more particularly to a faucet assembly with a valve which controls flow of water in a water mixing chamber to one of two outflow ports.\n2. Description of the Related Art\nReferring to FIGS. 1 and 2, a conventional faucet assembly is shown to include a faucet body 10 which has a water mixing chamber 103 for introducing hot and cold water from hot and cold water inflow passages 101,102. A primary valve 11 is disposed in the water mixing chamber 103, and extends in an axial direction to connect with a handle 12 so as to be rotatable to control the flow of the mixed water formed in the water mixing chamber 103 into an accommodating cavity 104 which is in fluid communication with the water mixing chamber 103. The accommodating cavity 104 is further in fluid communication with first and second outlets 105,106, such as a faucet spout and a showerhead. A flow direction controlling body 108 is disposed downstream of the water mixing chamber 103 and upstream of the second outlet 106. The flow direction controlling body 108 has an annular wall 107 which extends in the axial direction such that a secondary valve 13 is disposed sealingly and rotatably in the flow direction controlling body 108 by means of seal rings 131. A handle 14 has a spline connection with an outer portion 133 of the valve 13 to actuate rotation of the secondary valve 13. The valve 13 has a through hole 132 which extends radially and which is rotatable to communicate the accommodating cavity 104 with the second outlet 106 so as to permit the mixed water to flow out from the second outlet 106. When the through hole 132 is rotated so as not to register with the second outlet 106, the mixed water will flow out from the first outlet 105.\nHowever, a clearance arises between the outer surrounding wall of the valve 13 and the inner peripheral wall of the flow direction controlling body 22, thereby resulting in dripping at the second outlet 106.\nThe object of the present invention is to provide a faucet assembly which can overcome the aforementioned problem commonly associated with the prior art.\nAccording to this invention, the faucet assembly comprises a water mixing chamber which includes first front and rear inner walls facing towards and spaced apart from each other in an axial direction to define an accommodating space for receiving water. The first rear inner wall has hot and cold water inlets which are disposed to introduce hot and cold water flows into the accommodating space, and which can communicate with hot and cold water inflow passages, respectively. A transit chamber is in fluid communication with and downstream of the water mixing chamber. A valve assembly is disposed downstream of the mixing chamber and upstream of the transit chamber, and includes a primary valve with first and second valve portions. The first valve portion is disposed in the water mixing chamber, and is rotatable relative to the first rear inner wall about a first axis oriented in the axial direction to control the flows of hot and cold waters into the water mixing chamber via the hot and cold water inlets, respectively. The second valve portion is disposed between the water mixing chamber and the transit chamber to control the flow of mixed water formed in the water mixing chamber into the transit chamber. The valve assembly further includes a regulating member which is disposed in the water mixing chamber, and which is coupled to be rotatable with the first valve portion, and to be associated with the second valve portion. The regulating member extends in the axial direction and outwardly of the first front inner wall so as to be externally actuated to rotate the first valve portion and to control the second valve portion. A flow direction controlling body includes second front and rear inner walls which face toward and which are spaced apart from each other in the axial direction, and an inner tubular wall which surrounds a second axis oriented in the axial direction, which is interposed between, and which cooperates with the second front and rear inner walls to confine an accommodation cavity. The inner tubular wall has a communicating port which is disposed downstream of and which is communicated with the transit chamber, and which extends in a radial direction relative to the second axis to communicate the transit chamber with the accommodation cavity. The second rear inner wall has first and second outflow ports which are spaced apart from each other and which can respectively communicate with first and second outlets. A secondary valve includes a bottom wall with an annular periphery. The bottom wall is disposed to face toward and is in sliding contact with and rotatable relative to the second rear inner wall between first and second positions. The secondary valve further includes an annular wall which extends from the annular periphery in the axial direction towards the second front inner wall. The annular wall includes proximate and distal portions relative to the bottom wall. The distal portion is in sealing and sliding contact with the inner tubular wall when the bottom wall is rotated relative to the second rear inner wall. The proximate portion has a first passage extending in a radial direction relative to the second axis and communicating with the communicating port. The bottom wall has a second passage which extends therethrough in the axial direction to communicate with the first passage and which further extends angularly about the second axis. As such, when the bottom wall is rotated to the first position, the second passage communicates with the first outflow port. When the bottom wall is moved to the second position, the second passage communicates with the second outflow port. The secondary valve further includes a stem portion which is secured to the distal portion and which extends in the axial direction and outwardly of the second front inner wall so as to be actuated externally to rotate-the-bottom wall."}
-{"text": "There are a variety of devices the purpose of which is to evaporate the fuel components entering the intake in the liquid state by preheating the fuel-air mixture. In this manner, a better preparation of the fuel-air mixture is achieved, whereby the presence of constituents injurious to health in the exhaust may be diminished. Preheating of this type is advantageous especially at low load, since it prevents the carbureted fuel from recondensing in this operating mode, with resulting combustion of a leaner fuel-air mixture, thereby reducing the efficiency of the engine. At full load when the engine is hot, on the other hand, preheating is not desirable since losses of fill would result. To obtain controllable preheating, more or less effective depending on the operating condition of the engine, dampers have been proposed, arranged in the intake and/or exhaust line in order to deflect or divert the incoming fuel-air mixture and/or the exhaust flow from the heat transfer wall. Such dampers, however, have the disadvantage of interfering seriously with the direction of flow of the incoming fuel-air mixture by a deflection or partial constriction and displacement of the port area, and may lead to separation of the mixture, while in the exhaust line a resistance is created that undesirably increases the back pressure."}
-{"text": "The present invention relates to tumor necrosis factor receptor (TNF-R) related polypeptides (proteins) and their ligands, hereinafter referred to as TR1, TR2, TL2 and TL4. The invention also relates to inhibiting or activating the action of such polypeptides using agonists or antagonists by the screening methods described herein.\nMany biological actions are a response to certain stimuli and natural biological processes, and are controlled by factors, such as cytokines. These cytokines act through target cell receptors by engaging the receptor and producing an intracellular response.\nFor example, tumor necrosis factors (TNF) alpha and beta are cytokines which act through TNF receptors to regulate numerous biological processes, including protection against infection and induction of shock and inflammatory disease. The TNF molecules belong to the xe2x80x9cTNF-ligandxe2x80x9d superfamily, and act together with their receptors or counter-ligands, the xe2x80x9cTNF-receptorxe2x80x9d superfamily. So far, ten members of the TNF ligand superfamily have been identified and thirteen members of the TNF-receptor superfamily have been characterized.\nAmong the ligands there are included TNF-a, lymphotoxin-a (LT-a, also known as TNF-b), LT-b (found in complex heterotrimer LT-a2-b), FasL, CD40L, CD27L, CD30L, 4-1BBL, OX40L and TRAIL ((Wiley et al.Immunity 3: 673-682 (1995)) All but one of these (LTa) are expressed as type II membrane proteins. The superfamily of TNF receptors includes the p55TNF receptor, p75TNF receptor, TNF receptor-related protein, FAS antigen or APO-1, CD40, CD27, CD30, 4-1BB, OX40, low affinity p75, NGF-receptor(Meager, A., Biologicals, 22:291-295 (1994)).\nMany members of the TNF-ligand superfamily are expressed by activated T-cells, implying that they are necessary for T-cell interactions with other cell types which underlie cell ontogeny and functions. (Meager, A., supra).\nConsiderable insight into the essential functions of several members of the TNF receptor family has been gained from the identification and creation of mutants that abolish the expression of these proteins. For example, naturally occurring mutations in the FAS antigen and its ligand cause lymphoproliferative disease (Watanabe-Fukunaga, R., et al., Nature 356:314 (1992)), perhaps reflecting a failure of programmed cell death. Mutations of the CD40 ligand cause an X-linked immunodeficiency state characterized by high levels of immunoglubulin M and low levels of immunoglobulin G in plasma, indicating faulty T-cell-dependent B-cell activation (Allen, R. C. et al., Science 259:990 (1993)). Targeted mutations of the low affinity nerve growth factor receptor cause a disorder characterized by faulty sensory innovation of peripheral structures (Lee, K. F. et al, Cell 69:737 (1992)).\nTNF and LT-a are capable of binding to two TNF receptors (the 55- and 75-kd TNF receptors). A large number of biological effects elicited by TNF and LT-a, acting through their receptors, include hemorrhagic necrosis of transplanted tumors, cytotoxicity, a role in endotoxic shock, inflammation, immunoregulation, proliferation and anti-viral responses, as well as protection against the deleterious effects of ionizing radiation. TNF and LT-a are involved in the pathogenesis of a wide range of diseases, including endotoxic shock, cerebral malaria, tumors, autoimmuine disease, AIDS and graft-host rejection (Beutler, B. and Von Huffel, C., Science 264:667-668 (1994)). Mutations in the p55 Receptor cause increased susceptibility to microbial infection.\nMoreover, an about 80 amino acid domain near the C-terminus of TNFR1 (P55) and Fas was reported as the xe2x80x9cdeath domain,xe2x80x9d which is responsible for transducing signals for programmed cell death (Tartaglia et al., Cell 74:845 (1993)). Other regions of the TNF receptor intracellular domain are responsible for the activation of transcription through NF-kB (Cheng and Baltimore Genes and Development 10: 963-973 (1996)). More recent evidence has suggested that receptors may induce signals in cells expressing membrane bound TNF family ligand in a process known as xe2x80x9creverse signalingxe2x80x9d (Wiley et al., J. Immunol. 157:3635-3639.\nThe effects of TNF family ligands and TNF family receptors are varied and influence numerous functions, both normal and abnormal, in the biological processes of the mammalian system. There is a clear need, therefore, for identification and characterization of such receptors and ligands that influence biological activity, both normally and in disease states. In particular, there is a need to isolate and characterize novel members of the TNF receptor family.\nThis indicates that these Tumor necrosis factor receptors (TNF-R) have an established, proven history as therapeutic targets. Clearly there is a need for identification and characterization of further members of Tumor necrosis factor receptor (TNF-R) family which can play a role in preventing, ameliorating or correcting dysfunctions or diseases, including, but not limited to, chronic and acute inflammation, arthritis, septicemia, autoimmune diseases (eg inflammatory bowel disease, psoriasis), transplant rejection, graft vs. host disease, infection, stroke, ischemia, acute respiratory disease syndrome, restenosis, brain injury, AIDS, Bone diseases, cancer (eg lymphoproliferative disorders), atheroschlerosis, and Alzheimers disease.\nThe present invention relates to tumor necrosis factor receptor (TNF-R) related polypeptides and their ligands, hereinafter referred to as TR1, TR2, TL2 and TL4. The invention also relates to methods to identify agonists and antagonists of TR1, TR2, TL2 and TL4. The agonists and antagonists thus identified can be used to treat chronic and acute inflammation, arthritis, septicemia, autoimmune diseases (eg inflammatory bowel disease, psoriasis), transplant rejection, graft vs. host disease, infection, stroke, ischemia, acute respiratory disease syndrome, restenosis, brain injury, AIDS, bone diseases, cancer (eg lymphoproliferative disorders), atheroschlerosis, and Alzheimers disease, among others, caused by imbalance of TR1, TR2, TL2 or TL4.\nxe2x80x9cTR1 or TR1 polypeptide or TR1 proteinxe2x80x9d refers, among others, generally to a polypeptide having the amino acid sequence set forth in SEQ ID NO:1 as well as polypeptides comprising the amino acid sequence of SEQ ID NO: 1; and polypeptides comprising the amino acid sequence which have at least 70% identity to that of SEQ ID NO:1 over its entire length. Furthermore, TR1 also refers to a polypeptide which comprises a sequence which has 70% identity to a fragment of SEQ ID NO: 1. xe2x80x9cTR1 or TR1 polypeptide or TR1 proteinxe2x80x9d also includes derivatives, such as fusion proteins, of the above polypeptides, and some of these derivatives are further illustrated below.\nxe2x80x9cTR2 or TR2 polypeptide or TR2 proteinxe2x80x9d refers, among others, generally to a polypeptide having the amino acid sequence set forth in SEQ ID NO:2 as well as polypeptides comprising the amino acid sequence of SEQ ID NO: 2; and polypeptides comprising the amino acid sequence which have at least 70% identity to that of SEQ ID NO:2 over its entire length. Furthermore, TR2 also refers to a polypeptide which comprises a sequence which has 70% identity to a fragment of SEQ ID NO: 2. xe2x80x9cTR2 or TR2 polypeptide or TR2 proteinxe2x80x9d also includes derivatives, such as fusion proteins, of the above polypeptides, and some of these derivatives are further illustrated below.\nxe2x80x9cTL2 or TL2 polypeptide or TL2 proteinxe2x80x9d refers, among others, generally to a polypeptide having the amino acid sequence set forth in SEQ ID NO:3 as well as polypeptides comprising the amino acid sequence of SEQ ID NO: 3; and polypeptides comprising the amino acid sequence which have at least 70% identity to that of SEQ ID NO:3 over its entire length. Furthermore, TL2 also refers to a polypeptide which comprises a sequence which has 70% identity to a fragment of SEQ ID NO: 3. xe2x80x9cTL2 or TL2 polypeptide or TL2 proteinxe2x80x9d also includes derivatives, such as fusion proteins, of the above polypeptides, and some of these derivatives are further illustrated below.\nxe2x80x9cTL4 or TL4 polypeptide or TL4 proteinxe2x80x9d refers, among others, generally to a polypeptide having the amino acid sequence set forth in SEQ ID NO:4 as well as polypeptides comprising the amino acid sequence of SEQ ID NO: 4; and polypeptides comprising the amino acid sequence which have at least 70% identity to that of SEQ ID NO:4 over its entire length. Furthermore, TL4 also refers to a polypeptide which comprises a sequence which has 70% identity to a fragment of SEQ ID NO: 4. xe2x80x9cTL4 or TL4 polypeptide or TL4 proteinxe2x80x9d also includes derivatives, such as fusion proteins, of above polypeptides, and some of these derivatives are further illustrated below.\ncDNA encoding polypeptide of SEQ ID NO:4 is contained in SEQ ID NO:5.\nxe2x80x9cPolypeptidexe2x80x9d refers to any peptide or protein comprising two or more amino acids joined to each other by peptide bonds or modified peptide bonds, i.e., peptide isosteres. xe2x80x9cPolypeptidexe2x80x9d refers to both short chains, commonly referred to as peptides, oligopeptides or oligomers, and to longer chains, generally referred to as proteins. Polypeptides may contain amino acids other than the 20 gene-encoded amino acids. xe2x80x9cPolypeptidesxe2x80x9d include amino acid sequences modified either by natural processes, such as posttranslational processing, or by chemical modification techniques which are well known in the art. Such modifications are well described in basic texts and in more detailed monographs, as well as in a voluminous research literature. Modifications can occur anywhere in a polypeptide, including the peptide backbone, the amino acid side-chains and the amino or carboxyl termini. It will be appreciated that the same type of modification may be present in the same or varying degrees at several sites in a given polypeptide. Also, a given polypeptide may contain many types of modifications. Polypeptides may be branched as a result of ubiquitination, and they may be cyclic, with or without branching. Cyclic, branched and branched cyclic polypeptides may result from posttranslation natural processes or may be made by synthetic methods. Modifications include acetylation, acylation, ADP-ribosylation, amidation, covalent attachment of flavin, covalent attachment of a heme moiety, covalent attachment of a nucleotide or nucleotide derivative, covalent attachment of a lipid or lipid derivative, covalent attachment of phosphotidylinositol, cross-linking, cyclization, disulfide bond formation, demethylation, formation of covalent cross-links, formation of cystine, formation of pyroglutamate, formylation, gamma-carboxylation, glycosylation, GPI anchor formation, hydroxylation, iodination, methylation, myristoylation, oxidation, proteolytic processing, phosphorylation, prenylation, racemization, selenoylation, sulfation, transfer-RNA mediated addition of amino acids to proteins such as arginylation, and ubiquitination. See, for instance, PROTEINSxe2x80x94STRUCTURE AND MOLECULAR PROPERTIES, 2nd Ed., T. E. Creighton, W. H. Freeman and Company, New York, 1993 and Wold, F., Posttranslational Protein Modifications: Perspectives and Prospects, pgs. 1-12 in POSTTRANSLATIONAL COVALENT MODIFICATION OF PROTEINS, B. C. Johnson, Ed., Academic Press, New York, 1983; Seifter et al., xe2x80x9cAnalysis for protein modifications and nonprotein cofactorsxe2x80x9d, Meth Enzymol (1990) 182:626-646 and Rattan et al., xe2x80x9cProtein Synthesis: Posttranslational Modifications and Agingxe2x80x9d, Ann NY Acad Sci (1992) 663:48-62.\nxe2x80x9cIdentityxe2x80x9d is a measure of the identity of nucleotide sequences or amino acid sequences. In general, the sequences are aligned so that the highest order match is obtained. xe2x80x9cIdentityxe2x80x9d per se has an art-recognized meaning and can be calculated using published techniques. See, e.g.: (COMPUTATIONAL MOLECULAR BIOLOGY, Lesk, A. M., ed., Oxford University Press, New York, 1988; BIOCOMPUTING: INFORMATICS AND GENOME PROJECTS, Smith, D. W., ed., Academic Press, New York, 1993; COMPUTER ANALYSIS OF SEQUENCE DATA, PART I, Griffin, A. M., and Griffin, H. G., eds., Humana Press, New Jersey, 1994; SEQUENCE ANALYSIS IN MOLECULAR BIOLOGY, von Heinje, G., Academic Press, 1987; and SEQUENCE ANALYSIS PRIMER, Gribskov, M. and Devereux, J., eds., M Stockton Press, New York, 1991). While there exist a number of methods to measure identity between two polynucleotide or polypeptide sequences, the term xe2x80x9cidentityxe2x80x9d is well known to skilled artisans (Carillo, H., and Lipton, D., SIAM J Applied Math (1988) 48:1073). Methods commonly employed to determine identity or similarity between two sequences include, but are not limited to, those disclosed in Guide to Huge Computers, Martin J. Bishop, ed., Academic Press, San Diego, 1994, and Carillo, H., and Lipton, D., SIAM J Applied Math (1988) 48:1073. Methods to determine identity and similarity are codified in computer programs. Preferred computer program methods to determine identity and similarity between two sequences include, but are not limited to, GCG program package (Devereux, J., et al., Nucleic Acids Research (1984) 12(1):387), BLASTP, BLASTN, FASTA (Atschul, S. F. et al., J Molec Biol (1990) 215:403).\nAs an illustration, by a polypeptide having an amino acid sequence having at least, for example, 70% xe2x80x9cidentityxe2x80x9d to a reference amino acid sequence of SEQ ID NO:2 is intended that the amino acid sequence of the polypeptide is identical to the reference sequence except that the polypeptide sequence may include up to thirty amino acid alterations per each 100 amino acids of the reference amino acid of SEQ ID NO: 2. In other words, to obtain a polypeptide having an amino acid sequence at least 70% identical to a reference amino acid sequence, up to 30% of the amino acid residues in the reference sequence may be deleted or substituted with another amino acid, or a number of amino acids up to 30% of the total amino acid residues in the reference sequence may be inserted into the reference sequence. These alterations of the reference sequence may occur at the amino or carboxy terminal positions of the reference amino acid sequence or anywhere between those terminal positions, interspersed either individually among residues in the reference sequence or in one or more contiguous groups within the reference sequence.\nWe have now discovered that TL2 of SEQ ID NO: 3 (otherwise known as TRAIL or Apo-2L, (Wiley S R, et al., Immunity (6):673-682 (1995); Pitt et al., J. Biol. Chem. 271: 12687-12690 (1996)) is a ligand of TR1 of SEQ ID NO:1 (otherwise known as osteoprotegerin (OPG), W. S. Simonet, et al., Cell, Vol 89, pp 309-319, 1997). Further, we also discovered that TL4 of SEQ ID NO: 4 is a ligand of TR2 of SEQ ID NO: 2 (described by R. I. Montgomery et al., Cell, Vol 87, pp427-436, 1996; Kwon et al., J. Biol. Chem. 272: 14272-14276 (1997); Hsu et al., J. Biol. Chem. 272:13471-13474 (1997))). Thus, the TR1 and TR2 polypeptides of the present invention, and their respective ligands, TL2 and TL4, can be employed in a screening process for compounds which bind to the receptors, or to their ligands, and which activate (agonists) or inhibit activation of (antagonists) TR1 and TR2 receptor polypeptides of the present invention, or their respective ligands TL2 and TL4. Thus, polypeptides of the invention may be used to assess the binding of small molecule substrates, receptors and ligands in, for example, cells, cell-free preparations, chemical libraries, and natural product mixtures. These substrates, receptors and ligands may be natural substrates and ligands, or may be structural or functional mimetics. See Coligan et al., Current Protocols in Immunology 1(2):Chapter 5 (1991).\nTR1, TR2, TL2 and TL4 polypeptides are responsible for many biological functions, including many pathologies. Accordingly, it is desirous to find compounds and drugs which stimulate TR1, TR2, TL2 or TL4 on the one hand, and which can inhibit the function of TR1, TR2, TL2 or TL4 or remove TR1, TR2, TL2 or TL4 expressing cells on the other hand (also defined as antagonists). Antagonists for TR1, TR2, TL2, and TL4 (including agents which remove TR1, TR2, TL2 or TL4 expressing cells) may be employed for a variety of therapeutic and prophylactic purposes for such conditions as chronic and acute inflammation, arthritis, septicemia, autoimmune diseases (e.g. inflammatory bowel disease, psoriasis), transplant rejection, graft vs. host disease, infection, stroke, ischemia, acute respiratory disease syndrome, restenosis, brain injury, AIDS, Bone diseases, cancer (e.g. lymphoproliferative disorders), atheroschlerosis, and Alzheimers disease. Agonists can be employed for therapeutic and prophylactic purposes for such conditions responsive to activation of T cells and other components of the immune system, such as for treatment of cancer and AIDS. However, agonists can also be employed for inappropriate stimulation of T cells and other components of the immune system which leads to down modulation of immune activity with therapeutic or prophylactic application for conditions such, as chronic and acute inflammation, arthritis, septicemia, autoimmune diseases (e.g. inflammatory bowel disease, psoriasis), transplant rejection, graft vs. host disease, infection, stroke, ischemia, acute respiratory disease syndrome, restenosis, brain injury, bone diseases, atheroschlerosis, and Alzheimers disease.\nAntagonists may be identified using assays to detect compounds which inhibit binding of TL2 to TR1 (or TL4 to TR2) in either cell-free or cell based assays. Suitable cell-free assays may be readily determined by one of skill in the art. For example, an ELISA format may be used in which purified TR1 (or TR2), or a purified derivative of TR1 (or TR2), such as a fusion protein, containing the extracellular domain of TR1 (or TR2), is immobilized on a suitable surface, either directly or indirectly (e.g., via an antibody to TR1 (or TR2) or to the fused epitope or protein domain) and candidate compounds are identified by their ability to block binding of purified soluble, extracellular domain of TL2 to TR1 (or soluble, extracellular domain of TL4 to TR2). The binding of TL2 to TR1 (or TL4 to TR2) could be detected by using a label directly or indirectly associated with TL2 (or TL4). Suitable detection systems include the streptavidin horseradish peroxidase conjugate, or direct conjugation by a tag, e.g., fluorescein. Conversely, purified, soluble TL2 (or TL4) may be immobilized on a suitable surface, and candidate compounds identified by their ability to block binding of purified TR1 to TL2 (or TR2 to TL4). The binding of TR1 to TL2 (or TR2 to TL4) could be detected by using a label directly or indirectly associated with TR1 (or TR2). Many other assay formats are possible that use the TR1 (or TR2) protein and its ligands.\nSuitable cell based assays may be readily determined by one of skill in the art. In general, such screening procedures involve producing appropriate cells which express the receptor polypeptides (or ligands thereof) of the present invention on the surface thereof. Such cells include cells from mammals, yeast, Drosophila or E. coli. Cells expressing the receptor, such as TR1 or TR2, (or cell membrane containing the expressed receptor) are then contacted with a ligand, such as TL2 or TL4, or test compound to observe binding, or stimulation or inhibition of a functional response. The assays may simply test binding of a candidate compound wherein adherence to the cells bearing the receptor is detected by means of a label directly or indirectly associated with the candidate compound or in an assay involving competition with a labeled competitor, such as the ligands TL2 or TL4. Alternatively, cells expressing the ligand, such as TL2 or TL4, (or cell membrane containing the expressed ligand) are then contacted with a receptor, such as TR1 or TR2, or test compound to observe binding, or stimulation or inhibition of a functional response. Similarly, the assays may simply test binding of a candidate compound wherein adherence to the cells bearing the ligand is detected by means of a label directly or indirectly associated with the candidate compound or in an assay involving competition with a labeled competitor, such as the receptors TR1 and TR2. Further, these assays may test whether the candidate compound results in a signal generated by activation of the receptor (e.g TR1 or TR2) or its respective ligand (e.g. TL2 or TL4) using detection systems appropriate to the cells bearing the receptor or its ligand and fusion proteins thereof at their surfaces. Typical fusion partners include fusing the extracellular domain of the receptor or ligand with the intracellular tyrosine kinase domain of a second receptor. Inhibitors of activation are generally assayed in the presence of an agonist, such as the ligand TL2 or TL4 for cells expressing TR1 and TR2 receptors or receptor fusions respectively, or the receptor TR1 and TR2 with cells expressing TL2 and TL4 ligands and ligand fusions, and the effect on activation by the agonist by the presence of the candidate compound is observed. Standard methods for conducting such screening assays are well understood in the art.\nExamples of potential TR1 or TR2 antagonists include antibodies or, in some cases, proteins which are closely related to the ligand of the TR1 or TR2, e.g., a fragment of the respective ligand TL2 or TL4, or small molecules which bind to the receptor, or its ligand, but do not elicit a response, so that the activity of the receptor is prevented. Examples of potential TR1 or TR2 agonists include antibodies that bind to TR1 or TR2, its respective ligand, such as TL2 or TL4, or derivatives thereof, and small molecules that bind to TR1 or TR2. These agonists will elicit a response mimicking all or part of the response induced by contacting the native ligand.\nSince receptors may also invoke signals in cells expressing the membrane TL2 and TL4, these screens may also yield agonists which mimic the agonist activity of TR1 with membrane TL2 and TR2 with membrane TL4. Examples of potential TL2 or TL4 agonists include antibodies that bind to TL2 or TL4, its respective receptor, such as TR1 or TR2, or derivatives thereof, and small molecules that bind to TL2 or TL4. These agonists will elicit a response mimicking all or part of the response induced by contacting the native ligand. Alternatively, TR1 or TR2 may be expressed as a soluble protein, including versions which fuse all or part of TR1 or TR2 with a convenient partner peptide for which detection reagents are available, eg TR1-IgG or TR2-IgG fusions, and used in a solid state or solution phase binding assay. For example, the soluble TR1 or TR2 can be used to detect agonist or antagonist binding directly through changes that can be detected experimentally, eg surface plasmon resonance, nuclear magnetic resonance spectrometry, sedimentation, calorimetry. The soluble TR1 or TR2 can be used to detect agonist or antagonist binding indirectly by looking for competition of the candidate agonist or antagonist with a ligand, such as TL2 or TL4, whose binding can be detected. Ligand detection methods include antibody recognition, modification of the ligand via radioactive labeling, chemical modification (eg biotinylation), fusion to an epitope tag. Methods include ELISA based assays, immunoprecipitation and scintillation proximity.\nAssays similar to those described above using soluble or membrane bound TR1 or TR2 may also be used to identify and purify additional natural ligand(s) of TR1 or TR2. These ligands may be agonists or antagonists of the receptor.\nThus the invention relates to:\nI. A method for identifying agonists or antagonists to TR1 or TR2 comprising:\n(a) contacting a candidate compound with TR1 or TR2 in the presence of labeled or unlabeled ligand TL2 or TL4 respectively; and\n(b) assessing the ability of said candidate compound to compete with TL2 or TL4 binding to TR1 or TR2 respectively;\nII. The method of I in which TR1 or TR2 is on the surface of a host cell, on a cell membrane or on a solid support;\nIII. The method of II for identifying agonists which further includes determining whether the candidate compound affects a signal generated by TR1 or TR2 polypeptide at the surface of the cell, wherein a candidate compound which increases production of said signal is identified as an agonist;\nIV. An agonist identified by the method of I, II or III;\nV. The method of II for identifying antagonists which further includes determining whether the candidate compound affects a signal generated by TR1 or TR2 polypeptide at the surface of the cell, wherein a candidate compound which diminishes production of said signal is identified as an antagonist;\nVI. An antagonist identified by the method of I, II or V;\nVII. A method for identifying agonists or antagonists to TL2 or TL4 comprising:\n(a) contacting a candidate compound with TL2 or TL4 in the presence of labeled or unlabeled TR1 or TR2 respectively; and\n(b) assessing the ability of said candidate compound to compete with TR1 or TR2 binding to TL2 or TL4 respectively;\nVIII. The method of VII in which TL2 or TL4 is on the surface of a host cell, on a cell membrane or on a solid support;\nIX. The method of VIII for identifying agonists to TL2 or TL4 which includes determining whether the candidate compound affects a signal generated by TL2 or TL4 polypeptide at the surface of the cell, wherein a candidate compound which increases production of said signal is identified as an agonist;\nX. An agonist identified by the method of IX;\nXI. The method of VIII for identifying antagonists which further includes determining whether the candidate compound affects a signal generated by TL2 or TL4 polypeptide at the surface of the cell, wherein a candidate compound which diminishes production of said signal is identified as an antagonist; and\nXII. An antagonist identified by the method XI.\nThus in another aspect, the present invention relates to a screening kit for identifying agonists, antagonists, ligands, receptors, substrates, enzymes, etc. for TR1, TR2, TL2 or TL4 polypeptides; which comprises:\n(a) a TR1, TR2, TL2 or TL4 polypeptide, preferably that of SEQ ID NO:1, 2, 3 or 4;\n(b) a recombinant cell expressing a TR1, TR2, TL2 or TL4 polypeptide, preferably that of SEQ ID NO: 1, 2, 3or 4;\n(c) a cell membrane expressing a TR1, TR2, TL2 or TL4 polypeptide; preferably that of SEQ ID NO: 1,2,3 or4; or\n(d) antibody to a TR1, TR2, TL2 or TL4 polypeptide, preferably that of SEQ ID NO: 1, 2, 3 or 4. It will be appreciated that in any such kit, (a), (b), (c) or (d) may comprise a substantial component.\nProphylactic and Therapeutic Methods\nThis invention provides methods of treating an abnormal conditions related to both an excess of and insufficient amounts of TR1, TR2, TL2 or TL4 polypeptide activity.\nIf the activity of TR1, TR2, TL2 or TL4 polypeptide is in excess, several approaches are available. One approach comprises administering to a subject an inhibitor compound (antagonist) as hereinabove described along with a pharmaceutically acceptable carrier in an amount effective to inhibit activation by blocking binding of ligands to the TR1 or TR2 polypeptide, or by inhibiting a second signal, and thereby alleviating the abnormal condition.\nIn another approach, soluble forms of TR1 or TR2 polypeptides still capable of binding the ligand in competition with endogenous TR1 or TR2 polypeptide may be administered. Typical embodiments of such competitors comprise fragments of the TR1 or TR2 polypeptide.\nFor treating abnormal conditions related to an under-expression of TR1, TR2, TL2 or TL4 and its activity, several approaches are also available. One approach comprises administering to a subject a therapeutically effective amount of a compound which activates TR1, TR2, TL2 or TL4 polypeptide, i.e., an agonist as described above, in combination with a pharmaceutically acceptable carrier, to thereby alleviate the abnormal condition.\nFormulation and Administration\nAgonists and antagonist of TR1, TR2, TL2 or TL4 may be formulated in combination with a suitable pharmaceutical carrier. Such formulations comprise a therapeutically effective amount of the polypeptide or compound, and a pharmaceutically acceptable carrier or excipient. Such carriers include but are not limited to, saline, buffered saline, dextrose, water, glycerol, ethanol, and combinations thereof. Formulation should suit the mode of administration, and is well within the skill of the art. The invention further relates to pharmaceutical packs and kits comprising one or more containers filled with one or more of the ingredients of the aforementioned compositions of the invention.\nPolypeptides and other compounds of the present invention may be employed alone or in conjunction with other compounds, such as therapeutic compounds.\nPreferred forms of systemic administration of the pharmaceutical compositions include injection, typically by intravenous injection. Other injection routes, such as subcutaneous, intramuscular, or intraperitoneal, can be used. Alternative means for systemic administration include transmucosal and transdermal administration using penetrants such as bile salts or fusidic acids or other detergents. In addition, if properly formulated in enteric or encapsulated formulations, oral administration may also be possible. Administration of these compounds may also be topical and/or localized, in the form of salves, pastes, gels and the like.\nThe dosage range required depends on the choice of peptide or compound, the route of administration, the nature of the formulation, the nature of the subject\"\"s condition, and the judgment of the attending practitioner. Suitable dosages, however, are in the range of 0.1-100 xcexcg/kg of subject. Wide variations in the needed dosage, however, are to be expected in view of the variety of compounds available and the differing efficiencies of various routes of administration. For example, oral administration would be expected to require higher dosages than administration by intravenous injection. Variations in these dosage levels can be adjusted using standard empirical routines for optimization, as is well understood in the art."}
-{"text": "The present invention relates to apparatus for storing and dispensing particulate material and in particular to a toner or developer cartridge for electrostatographic printing machines.\nIn an electrostatographic reproducing apparatus commonly in use today, a photoconductive insulating member is typically charged to a uniform potential and thereafter exposed to a light image of an original document to be reproduced. The exposure discharges the photoconductive insulating surface in exposed or background areas and creates an electrostatic latent image on the member which corresponds to the image areas contained within the usual document. Subsequently, the electrostatic latent image on the photoconductive insulating surface is made visible by developing the image with developing powder referred to in the art as toner. Most development systems employ a developer material which comprises both charged carrier particles and charged toner particles which triboelectrically adhere to the carrier particles. During development the toner particles are attracted from the carrier particles by the charge pattern of the image areas in the photoconductive insulating area to form a powder image on the photoconductive area. This image may subsequently be transferred to a support surface such as copy paper to which it may be permanently affixed by heating or by the application of pressure. Following transfer of the toner image to a support surface, the photoconductive insulating member is cleaned of any residual toner that may remain thereon in preparation for the next imaging cycle.\nAs the toner particles are depleted from the developer material it is necessary to dispense additional toner particles into the developer mixture. Further, the carrier has a limited life due to a variety of problems occurring with continued use. For example, the carrier may become impacted with toner thereby altering its triboelectric properties or if the carrier is coated with the material to enhance triboelectric properties the coating may deteriorate with time. U.S. Pat. No. 4,614,165 to Folkins et al. describes a developing process involving the addition of both toner particles and carrier particles to the developer in the developer housing to insure that the usable life of the developer material in the chamber at any point in time is at least equal to the life of the photographic printing machine. To accommodate the addition of further carrier and toner material, waste or spent developer is removed from the developer housing when it seems to exceeds a predetermined quantity. Thus, in most commercial embodiments the concentration of toner particles within the developer mixture is maintained substantially constant. To achieve this particulate material containers for toner and/or developer which discharge the toner and/or developer into the development system have been used. In replenishing the toner particles in an electrostatographic printing machine, it is important to minimize any spillage which may result in contamination of otherwise uncontaminated areas. Further, the toner particles being very finely ground may become airborne carrying this contamination to other areas but immediately adjacent the development system. In addition, the spilled toner particles have a tendency to cling to an operator's hands or to the surrounding environment. Accordingly, there is a need to provide a toner and/or developer storing and dispensing device which may be used to periodically replenish the toner and/or developer in an automatic printing machine and which can avoid the above difficulties thereby overcoming a dirty and messy operation.\nWhile the desire to minimize contamination by the toner or developer and provide for \"white glove\" machine operator involvement during the replenishment activity of the toner or developer is of paramount importance it is also noted that the toner and/or the developer container must be capable of containing the toner and/or developer during periods of storage and/or shipment under conditions which may go from one extreme to another by way of temperature, pressure and/or impact. A common approach to sealing a toner container or cartridge has been to seal the dispensing opening with a breathable plastic film of polyvinyl chloride, for example, which has a pressure and heat actuated adhesive on one side which upon the application of heat and pressure forms a very good seal with the dispensing opening. When necessary for use, the dispensing opening can be opened and particulate materials emptied into the developer housing. A principle difficulty, however, associated with this type of sealing arrangement is the preciseness required for pressure, temperature and time in sealing the toner cartridge. For example, if a hot iron is left on to long, it may burn or melt through the plastic film. Accordingly, it has been found that the plastic film is generally too time temperature dependent. In addition, the adhesive is frequently applied unevenly and with the uneven application of the adhesive in trying to manually peel or remove the film from the opening the force that may have to applied may be uneven and an ineffective seal may be formed in areas where too little adhesive is applied. Further, once this plastic film seal is broken and upon exhaustion of the supply of toner and/or developer in the cartridge so that it must be removed and replenished with a full cartridge, there is really effectively no seal to protect the user from toner spillage and contamination upon loading the cartridge from the machine."}
-{"text": "1. Field of the Invention\nThis invention relates generally to automatic transmissions for trucks, and more specifically to automatic truck transmissions fitted with bypass valves for reducing dead time or delay in transmitting fluid control pressure signals when gears are changed manually.\n2. Brief Description of the Prior Art\nThe Detroit Diesel Allision Division of General Motors Corporation is the largest producer of automatic transmissions for truck and other heavy vehicles. An important product has been the MT 650 transmission with five forward speeds. This transmission was introduced in model year 1974. The construction, operation, and repair of this transmission are described in detail in Service Manual, Allison Transmissions, Automatic Models MT 640, 643, MT 650, 653 (Detroit Diesel Allision, Indianapolis, 1983), incorporated herein by reference. The MT 653 transmission is a modified version of the 650 transmission.\nIn the early production models of the MT 650 and 653 transmissions an extended time delay of from about five to forty seconds of clutch-up or engagement for reverse gear operation occured in some transmission applications. The delay was a result of the priority design of the hydraulic control system and occured when a shift was made from low gear to reverse gear. In the MT 650 and 653 transmissions, low gear, the lowest forward gear, is a very low gear which must be engaged manually, and first gear is the lowest forward gear which is shifted automatically. There are always two clutches engaged in the drive train for each gear. In first gear these are the first clutch and the forward clutch; in reverse gear these are the first clutch and the fourth clutch. In neutral, there is only one clutch applied, the first clutch. Thus, the shift from low gear to reverse gear involve releasing the forward and low clutches and applying the first and fourth clutches.\nIn an attempt to remedy the delay, the transmissions were modified by installing a bypass fitted with a check valve in between the reverse pressure line (fourth clutch circuit) and the first clutch apply pressure line (first clutch circuit). The check valve included a plastic ball valve member and was positioned to seal when the pressure in the first clutch apply pressure line was greater than the pressure in the reverse pressure line, and to open when the pressure in the reverse pressure line was greater than the pressure in the first clutch apply pressure line.\nThe plastic ball of the check valve would tend to dissolve i the transmission fluid. After a period of about three to five years, in particular in transmissions equipped with power-take off (\"PTO\") fittings for supplying hydraulic pressure to hydraulic cylinders used for sundry purposes in certain types of heavy vehicles, the ball would fail to seal when the transmission fluid pressure was raised in the first clutch pressure line.\nIn these transmissions, a forward clutch and the low clutch are simultaneously applied in operation of the first forward gear, and the first clutch and a fourth clutch are simultaneously applied in operation of reverse gear. However, if, when first gear is engaged and the second clutch pressure line is then pressurized, the check valve fails to seal, the reverse pressure line is also pressured, and the forward, first and fourth clutches are all applied simultaneously, resulting in severe damage to the transmission. This requires a time-consuming and expensive rebuilding of the failed transmission. If the ball were to fail when the transmission was in neutral, the first and fourth clutches would be unintentionally applied, and the vehicle would begin to move backwards, a potentially very dangerous situation.\nA number of different plastic materials for the ball, including nylon, were apparently substituted by the transmission manufacturer in attempts to solve this substantial problem, all without ultimate success.\nSince the time that installation of the plastic check valve in new model transmissions was initiated by the manufacturer, the manufacturer also altered the transmissions by replacing the input oil pressure pump with a \"wide gear\" version so as to increase the rate at which input oil pressure increases when the engine is accelerated Previously, the drive and driven gears of the oil pump had a width specified to be from 0.6835 inch to 0.6855 inch while the \"wide gear\" pump which replaced the earlier model has drive and driven gears which have a width specified to be from about 0.8180 inch thick to 0.8205 inch thick Two different driven gear widths and four different drive gear widths have been provided by the manufacturer for replacement, as slightly wider gears must be fit when a transmission is rebuilt and the oil pump cavity has become worn through circulation of oil suspending abrasive particles. The width of the drive gear is specified to be within plus or minus two and one half ten thousands of an inch and the with of the driven gear is specified to be within plus or minus five ten thousands of an inch.\nThe transmission manufacturer provided instructions for eliminating the delay in units in field service in \"S.I.L. 10-TR-80,\" Feb. 20, 1980, as well as in Instruction Sheet No. 230, used in connection with Service Kit Part Number 23041150, including a template and replacement separator plate for use in these transmissions. S.I.L. 10-TR-80 and Intruction Sheet No. 230 are each incorporated herein by reference.\nHeavy vehicles equipped with Allison MT 650 transmissions are widely used, and the rebuilding required after about three to five years service is extremely costly for owners of fleets of such vehicles, including utility companies. The costs include not only the direct cost of rebuilding the transmissions, but also the cost of maintaining larger fleets than would otherwise be necessary, the costs to customers for whom service is delayed because of transmission failure, and the like. Despite apparently repeated efforts to remedy the problem, the manufacturer has apparently been unsuccessful. There is a significant need for an improved transmission which simultaneously overcomes the problem of delay encountered when shifting from first to reverse gear, and the premature failure encountered after only about three to five years of service of transmissions having plastic ball check valves in between the first clutch apply pressure line and the reverse pressure line."}
-{"text": "It is known to provide chairs, especially stackable chairs as described in the Mauser office brochure entitled \"Mahrzweckstuhle Serie 570\" (Multipurpose Chair Series 570) of Mauser Office GmbH, Korbach, Germany, in which the chairs are described for auditoriums and the like which can be aligned in rows and connected together by coupling means between adjacent chairs. These chairs can be fitted with so-called lecture desks, i.e. a table carried by a mounting or holding mechanism enabling the table to be releasably attached to the chair so that the table is disposed just forwardly of the user and, ordinarily, on one side of the chair.\nThe lecture desk can be swung from a more or less horizontal position of use in which it is locked above the seating surface and forwardly of the user, into a substantially folded position alongside the seating service. The holding mechanism can, if desired, be attached to an armrest of the chair.\nFor assembling the chairs into rows for use in an auditorium or the like, the stackable chairs are removed from the stack, positioned side by side and connected together by appropriate coupling means. When the chairs are taken up from the rows, they can be stacked in a nested state so that space is conserved.\nThe coupling systems used for interconnecting the chairs in rows, are available in various configurations. For example, German patent 39 33 817 discloses a coupling system in which a stirrup is provided on the underside of the seating surface of one chair and hooks into a stirrup of the underside of an adjoining chair when these stirrups are extended toward one another. The free ends of the stirrups are formed with mutually-engaging hooks, A problem with this arrangement is that whether the chair is provided with or without an armrest, the lecture desk must be configured differently and different configurations of the couplings may have to be used.\nEarlier desk systems have not proved to be effective for chairs which could be provided with or without armrests, are adapted to be coupled together or to be free-standing, or are especially configured for assembly in rows."}
-{"text": "1. Field of the Invention\nThis invention relates to novel antibiotic compounds designated BU-3292T A and B. The compounds possess antimicrobial and antitumor activities.\n2. Description of the Prior Art\nElucidation of the structure of BU-3292T A revealed that it had a unique depsipeptide structure. This type of pentapeptide lactone has been proposed as a key intermediate in the biosynthesis of actinomycin (Katz. E: Antibiotics, Vol. II. Biosynthesis, p. 276, Springer-Verlay, Berlin, 1967), but none of such monomer type compound has so far been discovered in the microbial metabolites.\nActinomycin Z.sub.5 (Tetrahedron Letter, 2567, 1973) and Schering 1a and 1b (15th Interscience Conference on Antimicrobial Agents and Chemotherapy, Abstract No. 419) are reported to each contain one mole of 5-methylproline. BU-3292T B differs from these actinomycins in that it has two moles of 5-methylproline forming two identical peptide lactone side chains."}
-{"text": "As an electric current type of light-emitting device, organic light emitting diodes (OLEDs) have found wide application in high performance displays. With the increase in the size of the display, traditional passive matrix organic light-emitting diode (PMOLED) displays require shorter driving time for a single pixel, which necessitates an increased transient current and increased power consumption. Meanwhile, application of a large current results in a large voltage drop on the ITO wire and a high working voltage of the OLED, and in turn decreases its efficiency. Active matrix organic light-emitting diode (AMOLED) displays may address these issues elegantly by means of switch transistors scanning and inputting currents for OLEDs line by line.\nAmong the three types of driving approaches for AMOLED, which are digital driving, current driving and voltage driving, the voltage driving approach is similar to the traditional driving approach for active matrix liquid crystal displays (AMLCDs), i.e., providing by a driving chip (IC) a voltage signal that represents a grayscale, which voltage signal would be converted inside a sub-pixel into a current signal for driving a thin film transistor, so as to drive the OLED to exhibit a luminance grayscale. This approach is advantageous in that it is fast in driving speed, simple to implement, and appropriate for driving large-sized panels, and thus is extensively applied in the industry.\nHowever, in an AMOLED display of the voltage driving type, the current-luminance (I-L) conversion efficiency decreases as the OLED ages over time. Even if the currents are the same, the luminance displayed may be different because due to their different aging degrees the OLEDs are different in their conversion efficiencies. This leads to an issue that non-uniformity of the luminance is present in the images displayed on the display panel.\nIn order to improve the luminance uniformity, external compensation approaches are generally applied at present for the AMOLED display of the voltage driving type. Namely, each sub-pixel of the display panel is connected with a driving chip through a respective sense line corresponding one-to-one thereto, which driving chip detects the aging of the OLEDs in respective sub-pixels through the respective sense lines and then performs compensation of the sub-pixels in accordance with the detection. Nevertheless, in the above display panel, each sub-pixel is connected to a respective sense line, leading to increased wirings in the display panel, which is disadvantageous for fabrication of a high resolution display panel. Moreover, the number of the signal channels of the driving chip will be doubled, resulting in an increased area of the driving chip and a high cost."}
-{"text": "1. Field of the Invention\nThis invention relates to an improvement in a camera of the kind using a film provided with a magnetic storage part and having a magnetic head arranged in the camera to read information such as ISO data, etc. written in the magnetic storage part of the film and to write therein some information as necessary.\n2. Description of the Related Art\nCameras of the kind using a kind of film that has a magnetic storage part have recently been proposed, for example, as disclosed in U.S. Pat. No. 4,864,332. This camera is arranged to record in the magnetic storage part some photographic information such as a shutter time value, an aperture value, a date, a title, etc. and to read the recorded information as necessary out of the magnetic storage part by means of a magnetic head. However, practical attempts for affecting cameras of the kind having this magnetic head have encountered the following problems:\n1) Cameras of the kind disclosed in U.S. Pat. No. 4,841,319 etc. are arranged as follows: A film cartridge containing the whole roll of film completely therein including the leader part of the film is used. The film cartridge is inserted into a cartridge chamber, for example, from one side of the camera. The film is then pushed out of the cartridge by driving its shaft with a fork. With the cartridge of the kind used for a film having a magnetic storage part, an automatic film loading operation tends to become impossible, because the writing/reading magnetic head is merely arranged to come into contact with the film surface, the fore end part of the film tends to collide either with the magnetic head or with a film retaining pad and to be bent and warped when it is pushed out from the cartridge.\n2) When the camera is carried by the photographer, for example, by hanging on the neck while walking, vibrations are transmitted to the camera. Under such a condition, various accidents tend to happen including that the magnetic storage part of the film might be damaged by the movement of the magnetic head."}
-{"text": "Referring to FIG. 1, the nasal airway 1 comprises the two nasal cavities separated by the nasal septum, which airway 1 includes numerous ostia, such as the paranasal sinus ostia 3 and the tubal ostia 5, olfactory cells and is lined by the nasal mucosa. The nasal airway 1 can communicate with the nasopharynx 7, the oral cavity 9 and the lower airway 11, with the nasal airway 1 being in selective communication with the anterior region of the nasopharynx 7 and the oral cavity 9 by opening and closing of the oropharyngeal velum 13. The velum 13, which is often referred to as the soft palate, is illustrated in solid line in the closed position, as achieved by providing a certain positive pressure in the oral cavity 9, such as achieved on exhalation through the oral cavity 9, and in dashed line in the open position.\nThere are many nasal conditions which require treatment. One such condition is nasal inflammation, specifically rhinitis, which can be allergic or non-allergic and is often associated with infection and prevents normal nasal function. By way of example, allergic and non-allergic inflammation of the nasal airway can typically effect between 10 and 20% of the population, with nasal congestion of the erectile tissues of the nasal concha, lacrimation, secretion of watery mucus, sneezing and itching being the most common symptoms. As will be understood, nasal congestion impedes nasal breathing and promotes oral breathing, leading to snoring and sleep disturbance. Worryingly, the incidence of such allergic and non-allergic inflammatory diseases is increasing. Other nasal conditions include nasal polyps which arise from the paranasal sinuses, hypertrophic adenoids, secretory otitis media, sinus disease and reduced olfaction.\nIn the treatment of certain nasal conditions, the topical administration of medicaments is preferable, particularly where the nasal mucosa is the prime pathological pathway, such as in treating or relieving nasal congestion. Indeed, topical administration is advantageous in minimizing the possible side effects of systemic administration. Medicaments that are commonly topically delivered include decongestants, antihistamines, cromoglycates, steroids and antibiotics.\nThere are now an increasing number of adults and children who rely on pharmaceuticals to relieve symptoms associated with nasal conditions. At present, among the known anti-inflammatory pharmaceuticals, topical steroids have been shown to have an effect on nasal congestion. Topical decongestants have also been suggested for use in relieving nasal congestion. The treatment of hypertrophic adenoids and chronic secretory otitis media using topical decongestants, steroids and anti-microbial agents, although somewhat controversial, has also been proposed. Further, the topical administration of pharmaceuticals has been used to treat or at least relieve symptoms of inflammation in the anterior region of the nasopharynx, the paranasal sinuses and the auditory tubes.\nAside from the delivery of medicaments, the irrigation of the nasal mucosa with liquids, in particular saline solutions, is commonly practiced to remove particles and secretions, as well as to improve the mucociliary activity of the nasal mucosa. These solutions can be used in combination with active pharmaceuticals.\nFurthermore, medicaments are now increasingly systemically delivered through the nasal pathway, the nasal pathway offering a good administration route for the systemic delivery of pharmaceuticals, such as hormones, for example oxytocin, and antimigraine compositions, as the high blood flow and large surface area of the nasal mucosa advantageously provides for rapid systemic uptake.\nA variety of delivery systems have been developed to deliver substances to the nasal airways of subjects.\nConventionally, spray bottles have been used to deliver a medicament-containing liquid or an irrigating liquid to the nasal airways of subjects. However, the distribution of the delivered substance, in particular to the posterior region of the nasal airway, is less than ideal, especially in the cases of moderate and severe nasal obstruction. This poor distribution is often further exacerbated by a subject inhaling through the nasal airway during delivery, as is often prescribed, in an attempt to deliver the substance to the posterior region of the nasal airway. Indeed, an amount of the substance can be drawn into the lungs or swallowed in each delivery, which could be problematic in paediatric subjects if the medicament is a potent pharmaceutical, such as a steroid, which has to be administered frequently. In addition, the spray is frequently directed against the nasal septum which can undesirably lead to localized deposition. Further, the mechanical action of the delivery mechanism of the spray bottles can cause irritation and bleeding.\nGB-A-408856 discloses a delivery device, which, in one mode of use, allows for the inhalation of separate air flows entraining medicament into respective ones of the nasal cavities of a subject. This delivery device comprises a chamber containing a sponge saturated with medicament, a mouthpiece connected to the chamber and first and second nosepieces connected to the chamber. In one mode of use, the nosepieces are fitted into respective ones of the nostrils of a subject, and, on inhalation through the nosepieces, air flows entraining medicament are drawn into the lungs of the subject. In another mode of use, the mouthpiece is taken in the mouth of a subject, and, on inhalation through the mouthpiece, an air flow entraining medicament is drawn into the lungs of the subject.\nWO-A-98/53869 discloses a delivery device for delivering a powder containing a medicament to the nasal mucosa in one of the nasal cavities of a subject. This device comprises a tubular section which contains a metered dose of powdered medicament. In use, the ends of the tubular section are respectively located in the nostril of one of the nasal cavities and the mouth of a subject, and on exhalation by the subject through his or her mouth the exhaled air entrains the powdered medicament and delivers the same into the one nasal cavity, with the exhaled air backflowing out of the one nostril around the tubular section. In one embodiment the tubular section includes a flexible portion upstream of the dose of powdered medicament. The provision of this flexible portion allows the subject to close the tubular section at a point upstream of the medicament, such that, on release of the closed flexible portion during exhalation, a short explosive air flow entraining medicament is delivered into the one nasal cavity. In another embodiment the end of the tubular section located in the nostril can be shaped to act to locate the tubular section in a position in the nostril which allows for the deposition of the powdered medicament on the nasal mucosa.\nWhilst this delivery device is simple in construction, the operation of the device still does not provide for the effective delivery of substances, in particular one of a liquid or powder containing medicament, to the posterior region of the nasal airway, since medicament is delivered separately to each of the nasal cavities and the air flow into and out of each nasal cavity is through the same opening, namely the respective nostril, with the closed posterior region of the respective nasal cavity acting as a pressure reflecting surface which causes the exhaled air to backflow out of the one nostril before ever adequately reaching the posterior region of the respective nasal cavity. Further, in providing a short explosive burst of air flow into one of the nasal cavities, it is not possible to achieve a sustained and controlled bidirectional air flow through the nasal cavities which has been found necessary to deliver a substance effectively to the posterior region of the nasal airway.\nFor any substance to be delivered effectively to the nasal airway, it is highly desirable that the administration is efficient and simple. However, there can be problems in attempting to achieve this goal. In particular, the pathological changes observed with nasal inflammation make administration of substances, such as liquids or powders, tricky, particularly to the posterior region of the nasal airway and the posterior margins of the nasal structures. Indeed, as a consequence of the complex geometry of the narrow slit-like passages in the nasal airway, these passages become partially occluded when the nasal mucosa is inflamed and congested, making the distribution of topical pharmaceuticals to the nasal airway difficult."}
-{"text": "If a printing job is generated when an image forming apparatus is in a sleep state, the image forming apparatus returns from the sleep state and performs a printing preparation operation (hereinafter referred to as pre-run). A period from the start of the pre-run until printing is actually executed mainly depends on time for warming a fixing unit (hereinafter, heat accumulation time).\nIn the past, when the image forming apparatus returns from the sleep state, irrespective of what kind of a printing job is performed next, the image forming apparatus always warms the fixing unit until a fixed time elapses."}
-{"text": "Licensed wireless systems provide mobile wireless communications to individuals using wireless transceivers. Licensed wireless systems refer to public cellular telephone systems and/or Personal Communication Services (PCS) telephone systems. Wireless transceivers include cellular telephones, PCS telephones, wireless-enabled personal digital assistants, wireless modems, and the like.\nLicensed wireless systems utilize wireless signal frequencies that are licensed from governments. Large fees are paid for access to these frequencies. Expensive base station equipment is used to support communications on licensed frequencies. Base stations are typically installed approximately a mile apart from one another. As a result, the quality of service (voice quality and speed of data transfer) in wireless systems is considerably inferior to the quality of service afforded by landline (wired) connections. Thus, the user of a licensed wireless system pays relatively high fees for relatively low quality service.\nLandline (wired) connections are extensively deployed and generally perform at a lower cost with higher quality voice and higher speed data services. The problem with landline connections is that they constrain the mobility of a user. Traditionally, a physical connection to the landline was required. Currently, unlicensed wireless communication systems are deployed to increase the mobility of an individual using a landline. The mobility range associated with such systems is typically on the order of 100 meters. A common unlicensed wireless communication system includes a base station with a physical connection to a landline. The base station has a RF transceiver to facilitate communication with a wireless handset that is operative within a modest distance of the base station. Thus, this option provides higher quality services at a lower cost, but the services only extend a modest distance from the base station.\nThus, there are significant shortcomings associated with current landline systems and licensed wireless systems. For this reason, individuals commonly have one telephone number for landline communications and one telephone number for licensed wireless communications. This leads to additional expense and inconvenience for an individual. It would be highly desirable if an individual could utilize a single telephone number for both landline communications and licensed wireless communications. Ideally, such a system would allow an individual, through seamless handoffs between the two systems, to exploit the benefits of each system."}
-{"text": "1. Field of the Invention\nThe present invention relates to a solar collector for generating electricity and heat, and more particularly to a hybrid solar collector for converting solar energy of short wavelength into electric energy and converting solar energy of long wavelength into heat. The solar rays are separated into short wavelength and long wavelength. Separated short wavelength ray is focused on a solar cell for converting the energy of the short wavelength into electric energy, and separated long wavelength ray is focused on a pipe for heating a liquid inside the pipe by the energy of the long wavelength ray.\n2. Description of Related Art\nAs a solar energy collector for generating electric energy, a solar panel is well known. The solar panel is constituted of solar cells arranged in planar form. Further, a large-scale solar collector is also known. Such a solar collector is constituted of a large number of reflector mirrors distributed on the ground for converging solar rays at a converging point, and a solar cell arranged at the converging point.\nOn the other hand, a solar collector for generating heat is well known. Such a solar collector is constituted of reflector mirrors for converging solar rays, and a heat absorbing tube arranged at a converging point.\nMoreover, another solar collector for generating electricity and heat has been proposed. The hybrid solar collector utilizes a solar cell arranged at a solar rays converging point and a heat absorbing tube located on a rear surface of the solar cell."}
-{"text": "As is well known in the field of integrated circuit design, layout and fabrication, the manufacturing cost of a given integrated circuit is largely dependent upon the chip area required to implement desired functions. The chip area, in turn, is defined by the geometries and sizes of the active components such as gate electrodes in metal-oxide-semiconductor (MOS) technology, and diffused regions such as MOS source and drain regions and bipolar emitters and base regions. These geometries and sizes are often dependent upon the photolithographic resolution available for the particular manufacturing facility. The goal of photolithography in establishing the horizontal dimensions of the various devices and circuits is to create a pattern which meets design requirements as well as to correctly align the circuit pattern on the surface of the wafer. As line widths shrink smaller and smaller in submicron photolithography, the process to print lines and contact holes in photoresist becomes increasingly more difficult.\nWith circuit advancement to the ultra-large-scale integration (ULSI) levels, more and more layers are added to the surface of the wafer. These additional layers in turn create more steps on the wafer surface. The resolution of small image sizes in photolithography thus becomes more difficult over the additional steps due to the increased problem of depth of focus. Planarization techniques become increasingly more important to offset the effects of a varied topography.\nPlanarization techniques may be applied to both dielectric layers and conductors or semiconductors. Planarizing regions of a wafer may be done in degrees from smoothing a particular layer to reduce the severity of steep slopes to global planarization across the entire wafer regardless of the underlying topography. Of course, global planarization is ideal but is also more difficult to achieve particularly between widely isolated features. The variations in the thickness of the material to be etched, differences in adjacent materials being etched and the underlying topography all add to the degree of difficulty of achieving global planarization.\nInterconnect technology, creating connections between conducting regions, relies heavily on planarization techniques. In modern integrated circuits, the material of choice for upper-level conductive interconnection systems has been aluminum, including doped aluminum and aluminum alloys. Aluminum is an attractive material for integrated circuit metallization due to its high conductivity and low cost. The processing required to form aluminum metallization is also relatively easy, as it can readily be evaporated or sputtered onto the wafer. Aluminum is also able to form good ohmic contact to both p-type and n-type doped semiconductor material, such as silicon. In addition, aluminum is quite compatible with conventional semiconductor processes, such as that used to form bipolar and metal-oxide-semiconductor (MOS) devices, unlike other metals such as copper or gold which can diffuse into active regions and degrade device performance.\nCertain drawbacks do exist for aluminum-based metallization systems, however, particularly as geometries enter the submicron regime. A well-known limitation of aluminum is its poor step coverage, particularly for vertical or retrograde sidewalls of contact openings through insulating layers such as silicon dioxide, and especially for sputtered aluminum, due to the shadowing effect of steep contact walls. In addition, mechanical stress in the aluminum film can cause voids therein. Stress-induced voids and step coverage faults of sufficient size can each cause an open in a metal line or contact. Furthermore, since aluminum metallization is subject to electromigration, and since the rate of electromigration increases with current density through the film, necking or narrowing of an aluminum line due to such voids or poor steps locally increases the current density. As a result, the electromigration rate increases at narrowed locations of the film, greatly increasing the electromigration failure rate.\nA prior technique for addressing these limitations of aluminum metallization systems includes the use of refractory metal plugs, such as tungsten plugs, to fill contact openings in insulating layers. According to one example of this technique (for which many specific methods are well known in the art), after the opening of contacts through the insulating layer, a layer of tungsten is deposited by CVD over the wafer in such a manner as to conformally fill the contact opening, and is subsequently etched back to expose the surface of the insulating layer with the tungsten remaining in the contact opening. Alternatively, selective tungsten deposition has been used (the tungsten deposition on silicon but not on silicon dioxide) to fill contact openings. In either case, a subsequently deposited aluminum layer can readily make contact to the tungsten plug.\nWhile the tungsten plug technique has many advantages, including good step coverage in all contacts, compatibility with planarized processing, and tolerance of misalignment in the etch of overlying aluminum lines (since the aluminum can be etched selectively relative to the tungsten), the tungsten plug process adds complexity to the manufacturing flow. In addition, deposited tungsten is vulnerable to poor adhesion and high contact resistance, requiring the use of additional sputtered barrier films prior to the deposited tungsten, and the associated added process complexity therewith.\nAs circuit density and device performance requirements increase, more complex wiring or routing of interconnects between conductors is required. Vertical as well as horizontal interconnects are required as the number of metallization layers increase to meet the requirements of more complex circuits. Current generation devices may require up to 5 to 6 layers of metallization to meet the wiring needs of state-of-the-art devices such as logic devices and SRAMs. To make such devices manufacturable within the limitations of existing equipment, planarization at each stage of processing becomes even more critical.\nAccordingly, it is an object of the present invention to provide a method of forming planarized multilevel metallization regions to allow for complex routing of metal interconnections to achieve greater device density.\nIt is a further object of the present invention to form planarized regions containing metal interconnections and vertical metal pillars or vias between metal interconnections to support the multilevel metallization process.\nIt is yet a further object of the present invention to form planarized aluminum interconnections and vias in submicron geometries with adequate step coverage absent of stress-induced voids and electromigration problems.\nIt is a further object of the present invention to provide such a method which utilizes standard processing techniques.\nIt is yet a further object of the present invention to provide such a method which allows for closer spacing of interconnect lines enabling more devices to be formed in a smaller chip area.\nOther objects and advantages of the present invention will be apparent to those of ordinary skill in the art having reference to the following specification together with the drawings."}
-{"text": "1. Field of the Invention\nThis invention relates to an efficient method for neutralizing an alkoxylation catalyst. More particularly, this invention relates to a new use for partially spent ion exchange resins and an environmentally attractive method of neutralizing alkoxylation catalysts and removing potassium ions from polyol and alkoxylated alkyl phenol products. Example 2 demonstrates a reduction in potassium to less than 1 ppm.\nThis method is not only efficient and inexpensive, but also reduces problems inherent in disposal of a partially depleted acidic resin catalyst and reduces the impact on the environment of disposing of filter cakes which remain after using currently available methods to neutralize soluble base catalysts used in alkoxylation reactions.\n2. Description of Related Art\nOne of the most critical fields of study in industrial chemical processes is that of devising methods to remove impurities from products and effluent streams at low cost and, also, to avoid having to dispose of chemical compositions such as catalysts which are still partially active. Plants producing or using basic materials such as ammonia, organic amines, alkali salts or caustic have waste streams that often need to be neutralized. The manufacture of polyols and alkyl phenol alkoxylates presents such a problem.\nPolyols and alkyl phenol alkoxylates ar usually manufactured by the base catalyzed addition of propylene oxide or ethylene oxide to various initiators. The catalyst used is potassium hydroxide or another soluble base. The amount of potassium in the final product is critical in order to maintain control of downstream reactions. With one procedure, the catalyst is removed by neutralizing the base with a small excess of H.sub.2 SO.sub.4, absorbing the formed salt, (usually K.sub.2 SO.sub.4) with magnesium silicate, then filtering the product after heating to reduce the viscosity. The filter cake is then disposed of in a land fill.\nIn the past ion exchange resins have been used to produce polyols and, in fact, worked well. The disadvantage involved the fact that the solution used to regenerated the resin is considered a hazardous waste and is expensive to get rid of properly.\nThe use of ion exchange resins for removal of magnesium and noncarbonate minerals from water is known in the field of water treatment or water softening. The cost and frequency of regeneration are principal disadvantages. The ion-exchange material generally used in softening water are styrene-divinyl benzene copolymers. Kirk-Othmer Encyclopedia of Chemical Technology, 8, 70 and 24, 425 (1982).\nAnother aspect of the background for this invention is that the current transition to unleaded fuels in the U.S. has caused a demand for the addition of oxygenates into gasoline which has lead to the development of methyl-tert-butylether as a gasoline additive. The demand for MTBE has caused it to be the fastest growing chemical of the 80's and the demand will grow rapidly in the 90's. Chemical Business, January 1992, p.24. MTBE and other chemicals such as ethyl-tert-butylether (ETBE) and tert-amyl methylether (TAME) are made by the addition of alcohol to an olefin catalyzed by a sulfonic acid ion exchange resin.\nThe consumption of MTBE is currently about 180,000 barrels per day and could be as high as 670,000 barrels per day by the year 2000. Ibid, p. 25. Another reference indicates the demand would reach 1.2 million barrels per day by the year 2000. (Chemical Week, Nov. 20, 1991, p. 36.)\nDue to the demand for high productivity the ion exchange catalyst is used until reduced activity makes it more reasonable to exchange it for a fresh catalyst. This corresponds to a point of about 70-90%, and generally about 85%, of its original activity. This means there are enormous amounts of used catalyst that must be disposed of. Disposal of this partially spent catalyst gives rise to environmental concerns and often requires some sort of permit. It would be extremely efficient from a commercial and environmental viewpoint to accomplish a desirable goal using this partially spent catalyst, reduce the remaining activity of the resin so that disposal of it is not as objectionable and eliminate the production of a filter cake which requires disposal.\nIt would be a distinct advance in the art if an inexpensive means were available for neutralizing an alkoxylation catalyst and reducing impermissible levels of potassium in polyol and alkoxylated alkyl phenol products. It would be advantageous if this could be accomplished with a substance which to this point has had to be discarded while activity remains due to purity requirements of the primary product, i.e. MTBE, ETBE or TAME. To remove potassium from polyols and alkoxylated alkyl phenols with a partially spent catalyst, that normally entails disposal problems, and eliminate the neutralization of potassium hydroxide and the resulting filter cake would be quite efficient in every respect."}
-{"text": "This invention relates generally to marine propulsion devices, such as outboard motors, and, more particularly, to a drive shaft assembly for use in such a marine propulsion device.\nDrive shafts in conventional outboard motors are typically provided with one or more bearing races adjacent their lower ends and are exposed to corrosive water and exhaust gases at their upper ends. To provide improved lower end wear characteristics in conjunction with improved upper end corrosion resistance, prior drive shafts have been manufactured by inertially welding a lower drive shaft member of carbon steel to an upper drive shaft member of stainless steel. However, inertia welding is disadvantageous in that it requires costly secondary operations such as machining, heat treating and straightening.\nAttention is directed to the following U.S. Pat. No.:\nZackrisson 4,527,978 July 9, 1985 PA1 Federmann et al. 4,358,284 Nov. 9, 1982 PA1 Loyd et al. 4,272,971 June 16, 1981 PA1 Lakiza et al. 4,127,080 Nov. 28, 1978 PA1 Liaaen 4,112,866 Sept. 12, 1978 PA1 Campbell et al. 3,209,720 Oct. 5, 1965 PA1 Toelke 3,175,374 March 30, 1965 PA1 Shipley et al. 3,138,010 June 23, 1964 PA1 Willis, Jr., et al. 3,083,680 April 2, 1963"}
-{"text": "The invention relates to the field of cooling electronic devices and, in particular, to using circulating fluids to cool microprocessors, graphics processors, and other computer components.\nMicroprocessor dies typically used in personal computers are packaged in ceramic packages that have a lower surface provided with a large number of electrical contacts (e.g., pins) for connection to a socket mounted to a circuit board of a personal computer and an upper surface for thermal coupling to a heat sink. In the following description, a die and its package are referred to collectively as a microprocessor.\nElevation views of typical designs for heat sinks suggested by Intel Corporation for its Pentium\u00ae III microprocessor are shown in FIGS. 1A and 1B.\nIn FIG. 1A, a passive heat sink indicated generally by reference numeral 110 is shown. The passive heat sink 110 comprises a thermal plate 112 from the upper surface of which a number of fins, one of which is indicated by reference numeral 114, protrude perpendicularly. The passive heat sink 110 is shown in FIG. 1A installed upon a microprocessor generally indicated by reference numeral 118. The microprocessor 118 is comprised of a die 116 and a package 120. The die 116 protrudes from the upper surface of the package 120. The lower surface of the package 120 is plugged into a socket 122, which is in turn mounted on a circuit board (not shown). The passive heat sink 110 is installed by bringing the lower surface of the thermal plate 112 into contact with the exposed surface of the die 116. When installed and operated as recommended by the manufacturer, ambient airflow passes between the fins in the direction shown by an arrow 124 in FIG. 1A.\nIn FIG. 1B, an active heat sink, indicated generally by reference numeral 126, is shown. The active heat sink 126 comprises a thermal plate 128 from the upper surface of which a number of fins 130 protrude perpendicularly. A fan 132 is mounted above the fins 130. The active heat sink 126 is shown in FIG. 1B installed upon a microprocessor, generally indicated by reference numeral 136, which is comprised of a die 134 and a package 138. The die 134 protrudes from the upper surface of the package 138. The lower surface of the package 138 is plugged into a socket 140, which is in turn mounted on a circuit board (not shown). The active heat sink 126 is installed by bringing the lower surface of the thermal plate 128 into contact with the exposed surface of the die 134. When installed and operated as recommended by the manufacturer, ambient air is forced between the fins 130 in the direction shown by an arrow 142 in FIG. 1B.\nA difficulty with the cooling provided by the heat sinks shown in FIGS. 1A and 1B is that at best the temperature of the thermal plates 112, 128 can only approach the ambient air temperature. If the microprocessor 118, 136 is operated at a high enough frequency, the die 116, 134 can become so hot that it is difficult to maintain a safe operating temperature at the die 116, 134 using air cooling in the manner shown in FIGS. 1A and 1B.\nLiquid cooling, which is inherently more efficient due to the greater heat capacity of liquids, has been proposed for situations in which air cooling in the manner illustrated in FIGS. 1A and 1B is inadequate. In a typical liquid cooling system, such as that illustrated in FIG. 1C, a heat conductive block 144 having internal passages or a cavity (not shown) replaces the thermal plate 128 in FIG. 1B. The block 144 has an inlet and an outlet, one of which is visible and indicated by reference numeral 146 in FIG. 1C. Liquid is pumped into the block 144 through the inlet and passes out of the block 144 through the outlet to a radiator or chiller (not shown) located at some distance from the block 144. The block 144 is shown in FIG. 1C installed upon a microprocessor generally indicated by reference numeral 148, which is comprised of a die 150 and a package 152. The die 150 protrudes from the upper surface of the package 152. The lower surface of the package 152 is plugged into a socket 154, which is in turn mounted on a circuit board (not shown). The block 144 is installed by bringing its lower surface into contact with the exposed surface of the die 150.\nIn all liquid cooling systems known to the inventor, only a small portion of the lower surface of the block 144 comes into contact with the die 150. Since the die 150 protrudes above the upper surface of the package 152, a gap 156 remains between the upper surface of the package 152 and the block 144. If the gap 156 is not filled with insulation and sealed, convective and radiative heat transfer from the package 152 to the block 144 may occur. This will have no serious consequences so long as the block 144 is not cooled below the dew point of the air in the gap 156. If the liquid pumped through block 144 is only cooled by a radiator, then that liquid and consequently the block 144, can only approach the ambient air temperature. However, if a chiller is used to cool the liquid, then the temperature of the block 144 can decrease below the ambient air temperature, which may allow condensation to form on the package 152 or the block 144. Such condensation, if not removed, can cause electrical shorts, which may possibly destroy the microprocessor 148.\nCurrent solutions to the condensation problem referred to above include (1) controlling the chiller so that the temperature of the block 144 does not decrease below the dew point of the air in the gap 156 or (2) providing sufficient insulation and sealing material to prevent condensation from forming or to at least prevent any condensation that does form from reaching critical portions of the microprocessor 148 or surrounding circuit elements. Placing a lower limit on the temperature of the chiller limits the amount of heat that can effectively be removed from the microprocessor 148 without using bulky components. Further, the operating temperature of the microprocessor 148 can only approach the temperature of the block 144; operation at lower temperatures may be desirable in many circumstances. Alternatively, if insulation and sealing is used, trained technicians must do the installation properly if the installation is to be effective. If the insulation or seals fail, condensation can occur and cause catastrophic failure of the personal computer. A simpler, more reliable solution to the condensation problem is needed."}
-{"text": "1. Field of the Invention\nThis invention relates to fluid flow devices, particularly of that class of fluid flow nozzles with laminar discharge, and further to the field of illuminated fluid nozzles.\n2. Description of the Prior Art\nIn U.S. Pat. No. 5,160,086 granted on Nov. 3, 1992, a laminar flow fluid nozzle is provided for use in decorative water fountains, and also for industrial applications. Initial fluid flow is provided through a double-walled, bladder like fluid supply hose (32).\nThe inflated double hose isolates fluid flow from the pump and other vibrations which would otherwise be transmitted by means of a rigid hose. Furthmore, the double hose functions to absorb or accumulate small pressure variations known as \"pump noise\" which are usually present in the input stream.\nThe fluid next flows into a fluid chamber over a baffle (22). Inside pockets of air (18) cushion remaining pressure variations of the fluid stream.\nFluid next flows through a diffuser (20) which provides a purality of parallel fluid baths to dampen remaining major currents by lowering the fluid velocity, and thus the Reynolds number to obtain laminar fluid flow.\nA knife edged orifice (12) results in a narrow fluid stream (14). However, orifice (12) may be offset with respect to the center of enclosure (11).\nA light transmitting means (24) is aligned with orifice (12) by focusing means (28.\nA spring system (30) may be provided to absorb ambiant vibrations."}
-{"text": "1. Field of the Invention\nThe present invention relates to a power transmission apparatus using gears. More particularly, it relates to a power transmission apparatus having a first shaft and a second shaft each supported by an aluminum housing, and an iron gear and a synthetic resin gear which are mounted on to said two shafts respectively and are engaged with each other.\n2. Description of the Prior Art\nPreviously known materials used for a gear, a power transmission element, include those which comprise a synthetic resin as the main component (see, for example, Japanese Utility Model Application Kokai (Laid-Open) No. 70,702/85). Synthetic resin gears of this kind are used, for example, for camshaft gears of engines. An example of such uses will be described below with reference to the accompanying figure illustrating a working example of that utility model.\nIn the FIGURE, 1 indicates the crankshaft of an engine (a first shaft), 2 indicates the camshaft (a second shaft), 4a indicates a crankcase (housing) and 4b indicates the crankcase cover (housing). The crankshaft 1 and the camshaft 2 are arranged in parallel with each other and are supported by the crankcase 4a and the crankcase cover 4b. The crankshaft gear 5 made of iron press-fitted to the crankshaft 1 and the camshaft gear 7 made of synthetic resin fixed to the camshaft 2 with rivets 6 engage with each other and transmit the torque of the crankshaft 1 to the camshaft 2. In the prior art, as the material for the camshaft gear 7, synthetic resins containing with glass fibers have been used to improve the mechanical properties of the camshaft gear 7.\nSince the camshaft gear 7 is made of a synthetic resin, it is abraded by meshing with the crankshaft gear 5 made of iron even though the resin contains glass fibers, causing the glass fibers to fall off into lubricating oil. Since the glass fibers which fall off are minute, they get into, for example, such sliding parts as the bearing parts 8 and 9 of the crankshaft 1 and the camshaft 2. Since the glass fiber is very hard, the above-mentioned phenomenon brings about the disadvantage of abrasion of the sliding parts. Further, there is another disadvantage in that the glass fibers come up to the surface of the gear teeth on the side of the camshaft gear 7 made of resin, whereby the abrasion of the opposing crankshaft gear 5 of iron is increased.\nAs a means for solving such problems, it is possible to secure the necessary mechanical properties of the gear, without including glass fibers in the synthetic resin, by increasing the face width, changing to a resin with more enhanced heat resistance, or adding a substance capable of enhancing sliding such as molybdenum. However, the elimination of glass fibers increases the coefficient of linear expansion of the resulting material and markedly increases the change of backlash with temperature. Thus, the coefficient of linear expansion is 2.3.times.10.sup.-5 /deg C for aluminum, 1.2.times.10.sup.-5 /deg C for iron and 9.times.10.sup.-5 /deg C for the synthetic resin. Accordingly, with respect to the distance D between the crankshaft 1 and the crankshaft 2, the sum of the magnitude of expansion (or contraction) of the two gears 5 and 7 is larger than the magnitude of expansion (or contraction) of the crankcase 4, so that the backlash between the two gears 5 and 7 will differ greatly from the set value with changes in temperature.\nIf the change of backlash with temperature is large, when, for example, the backlash is so set as to be optimum at high temperatures, the backlash will become excessively large at low temperatures, causing emission of noise. On the contrary, when the backlash is so set as to be optimal at low temperatures, the backlash will become small at high temperatures, causing emission of noise and severe abrasion.\nOn the other hand, phenolic resins and the like containing cloth or wood chips have little of the above-described problems with respect to both contamination and backlash. However, generally they require the process of gear cutting, which markedly lowers the production efficiency."}
-{"text": "The present disclosure relates to an image reading device, an image forming apparatus including the same, and an image reading method.\nGenerally, there is known a light source switching type image reading device that obtains color image data by sequentially lighting light sources (LED, or the like) of three colors corresponding to R (red), G (green), and B (blue) and by reading, at that time, light data received by a line sensor as R data, G data, and B data, respectively.\nIn the light source switching type image reading device, since image data of respective colors are read while a document and the line sensor are relatively moved in a sub-scanning direction at a constant speed, the reading positions of the respective colors in the document do not coincide and are shifted in the sub-scanning direction. As a result, for example, each color is sometimes different from that in the document at, for example, a boundary portion of the color. Under such circumstances, there is known an image reading device that is able to correct such a color shift by correcting the image data of the respective colors that are read."}
-{"text": "1. Field of the Invention\nThe present invention relates to drainage for land areas. More particularly, the present invention relates to a system for draining land areas, such as golf courses and farms through the collection of water into a principal catch basin having permeable walls, and subsequently siphoning the water collected from the catch basin to a distant exit point such as a lake, canal, ditch or the like. There could be further included a series of secondary catch basins draining into the principal permeable catch basin, so that an even larger area may be drained through siphoning from a single control system.\n2. General Background\nIn the present state of the art, drainage water such as water from heavy rains or the like may be collected either by a catch basin or a sub-surface drainage system that collects water through seepage. At the present time, a catch basin, which consist of a substantially box-like chamber is set within the ground and positioned at a low point where water would flow under gravity and be collected within the catch basin. At this point the water is exited through a solid pipe at a one percent slope whenever possible, so that the water then exits into some type of collection system such as a lake, pond or creek. However, often times there is insufficient fall from the point of the collection system to the exit point to either obtain a 1% grade or any positive grade at all. The only alternatives, therefore, in the past, have been to either not attempt to do the drainage at all or to attempt to build a gravel sump, dry well, or sump pump. The concept behind these systems relate to the fact that the gravel sump will provide some temporary storage capacity for the water moving out of the collection system and maintain it until it can slowly seep back into the soil profile. Often times there will be a dry well built that will extend deep enough to reach a permeable strata below the less permeable layers that are near the surface. The system has two main problems, i.e., (a) the sumps have very little capacity, and once the storage capacity is filled, the collection system will back up and stop functioning; and (b) for a system to extend deep enough into the ground to hit a sand strata, it can be a very expensive system to install, depending on the depth of the sand strata. The Environmental Protection Agency is now concerned that bringing surface water down to these strata may contaminate the water table.\nThe other alternative is to build a sump pump which operates a pump off of a float switch and pumps the water out every time the level of the water activates the float switch. The problem with these systems is that they require a large amount of maintenance, are expensive, and a pump has to be located at each sump.\nThere also is a problem with the existing type of catch basin that is presently installed. These systems operate very well when they are surrounded by an impermeable surface, such as asphalt or concrete for street or parking lot drainage. However, in land drainage such as farms and golf courses, these systems are surrounded by permeable areas of soil. The problem is that before the water can enter the catch basin system, many times a large amount of water is trapped into the permeable soil profile surrounding the catch basin itself. Once that water is trapped into the soil profile, it cannot enter through the solid walls constructed of concrete or brick, and stays in the soil profile, thus creating a wet soggy saturated area. The other problem with existing catch basins is that they are designed only to stop the largest particles from entering the pipe that is used to drain the water out of the catch basin. Therefore, many contaminants such as small rocks and branches become lodged into the piping system. To prolong the time period before clogging occurs in drainage pipes, it has been a common practice to oversize the pipe in relation to the amount of water that is necessary to be drained. This oversizing of pipe to prevent clogging increases the expense of materials for the piping, as well as the installation costs because larger amounts of dirt must be excavated to install the pipe. Also, there is no system provided in existing drainage systems for allowing the escape of air trapped within the soil profile. This air reaches an equilibrium with the infiltrating water and stops the infiltration of the water into the soil profile. If this air were allowed to escape, a larger portion of the profile could absorb the infiltrating rainwater.\nThere have been several patents which address the drainage of water in a system. The most pertinent are listed as follows, including a narrative of each submitted in applicants Prior Art Statement, filed herewith:\n______________________________________ U.S. Pat. No. INVENTOR TITLE ______________________________________ 4,299,697 Curati, Jr. \"A Liquid Containment And Storage System For Rail- road Tracks\" 1,734,777 Pike \"System Of Draining\" 462,864 Hershberger \"Catch Basin For Sewers\" 2,432,203 Miller \"Catch Basin\" 4,472,086 Leach \"Geotextile Fabric Construction\" ______________________________________"}
-{"text": "The present invention provides compounds, pharmaceutical compositions containing one or more of those compounds or their pharmaceutically acceptable salts, which are effective in inhibiting the binding of various chemokines, such as MIP-1\u03b1, leukotactin, MPIF-1 and RANTES, to the CCR1 receptor. As antagonists or modulators for the CCR1 receptor, the compounds and compositions have utility in treating inflammatory and immune disorder conditions and diseases.\nHuman health depends on the body's ability to detect and destroy foreign pathogens that might otherwise take valuable resources from the individual and/or induce illness. The immune system, which comprises leukocytes (white blood cells (WBCs): T and B lymphocytes, monocytes, macrophages granulocytes, NK cell, mast cells, dendritic cell, and immune derived cells (for example, osteoclasts)), lymphoid tissues and lymphoid vessels, is the body's defense system. To combat infection, white blood cells circulate throughout the body to detect pathogens. Once a pathogen is detected, innate immune cells and cytotoxic T cells in particular are recruited to the infection site to destroy the pathogen. Chemokines act as molecular beacons for the recruitment and activation of immune cells, such as lymphocytes, monocytes and granulocytes, identifying sites where pathogens exist.\nDespite the immune system's regulation of pathogens, certain inappropriate chemokine signaling can develop and has been attributed to triggering or sustaining inflammatory disorders, such as rheumatoid arthritis, multiple sclerosis and others. For example, in rheumatoid arthritis, unregulated chemokine accumulation in bone joints attracts and activates infiltrating macrophages and T-cells. The activities of these cells induce synovial cell proliferation that leads, at least in part, to inflammation and eventual bone and cartilage loss (see, DeVries, M. E., et al., Semin Immunol 11(2):95-104 (1999)). A hallmark of some demyelinating diseases such as multiple sclerosis is the chemokine-mediated monocyte/macrophage and T cell recruitment to the central nervous system (see, Kennedy, et al., J. Clin. Immunol. 19(5):273-279 (1999)). Chemokine recruitment of destructive WBCs to transplants has been implicated in their subsequent rejection. See, DeVries, M. E., et al., ibid. Because chemokines play pivotal roles in inflammation and lymphocyte development, the ability to specifically manipulate their activity has enormous impact on ameliorating and halting diseases that currently have no satisfactory treatment. In addition, transplant rejection may be minimized without the generalized and complicating effects of costly immunosuppressive pharmaceuticals.\nChemokines, a group of greater than 40 small peptides (7-10 kD), ligate receptors expressed primarily on WBCs or immune derived cells, and signal through G-protein-coupled signaling cascades to mediate their chemoattractant and chemostimulant functions. Receptors may bind more than one ligand; for example, the receptor CCR1 ligates RANTES (regulated on activation normal Tcell expressed), MIP-1\u03b1 (macrophage inflammatory protein), MPIF-1/CK\u03b28, and Leukotactin chemokines (among others with lesser affinities). To date, 24 chemokine receptors are known. The sheer number of chemokines, multiple ligand binding receptors, and different receptor profiles on immune cells allow for tightly controlled and specific immune responses. See, Rossi, et al., Ann. Rev. Immunol. 18(1):217-242 (2000). Chemokine activity can be controlled through the modulation of their corresponding receptors, treating related inflammatory and immunological diseases and enabling organ and tissue transplants.\nThe receptor CCR1 and its chemokine ligands, including, for example MIP-1\u03b1, MPIF-1/CK\u03b28, leukotactin and RANTES, represent significant therapeutic targets (see Saeki, et al., Current Pharmaceutical Design 9:1201-1208 (2003)) since they have been implicated in rheumatoid arthritis, transplant rejection (see, DeVries, M. E., et al., ibid.), and multiple sclerosis (see, Fischer, et al., J Neuroimmunol. 110(1-2):195-208 (2000); Izikson, et al., J. Exp. Med. 192(7):1075-1080 (2000); and Rottman, et al., Eur. J. Immunol. 30(8):2372-2377 (2000). In fact, function-blocking antibodies, modified chemokine receptor ligands and small organic compounds have been discovered, some of which have been successfully demonstrated to prevent or treat some chemokine-mediated diseases (reviewed in Rossi, et al., ibid.). Notably, in an experimental model of rheumatoid arthritis, disease development is diminished when a signaling-blocking, modified-RANTES ligand is administered (see Plater-Zyberk, et al., Immunol Lett. 57(1-3):117-120 (1997)). While function-blocking antibody and small peptide therapies are promising, they suffer from the perils of degradation, extremely short half-lives once administered, and prohibitive expense to develop and manufacture, characteristic of most proteins. Small organic compounds are preferable since they often have longer half lives in vivo, require fewer doses to be effective, can often be administered-orally, and are consequently less expensive. Some organic antagonists of CCR1 have been previously described (see, Hesselgesser, et al., J. Biol. Chem. 273(25): 15687-15692 (1998); Ng, et al., J. Med. Chem. 42(22):4680-4694 (1999); Liang, et al., J. Biol. Chem. 275(25):19000-19008 (2000); and Liang, et al., Eur. J. Pharmacol. 389(1):41-49 (2000)). In view of the effectiveness demonstrated for treatment of disease in animal models (see, Liang, et al., J. Biol. Chem. 275(25):19000-19008 (2000)), the search has continued to identify additional compounds that can be used in the treatment of diseases mediated by CCR1 signaling."}
-{"text": "The production of polyesters and copolyesters of dicarboxylic acids and aliphatic glycols has been carried out commercially for many decades. Among the earliest disclosures relating to the production of polyesters and copolyesters is the disclosure in U.S. Pat. No. 2,465,319, issued Mar. 22, 1949. Since this disclosure many variations have been made in the process and many catalysts have been discovered and patented. On Dec. 8, 1970, there issued U.S. Pat. No. 3,546,179, which is directed to the use of compounds containing both silicon and phosphorus atoms a compounds as a catalyst for the production of such polyesters and copolyesters.\nThe use of metal halides and certain silicon compounds as catalysts in the polycondensation of dicarboxylic acids and aliphatic glycols is disclosed in U.S. Pat. Nos. 4,143,057; 4,254,241; and Re-30,554. The metal halides employed in the examples therein are necessarily employed with solvent present. In addition these compounds tend to be hydrolytically unstable and, thus, require careful handling during their use. Thus, although these catalysts are highly advantageous it would be highly desirable to have a catalyst that may be of use without a solvent and which is hydrolytically stable."}
-{"text": "In the manufacture of microelectronic devices, integrated circuits are formed on a semiconductor substrate, which generally includes silicon or other materials such as gallium arsenide, and a plurality of interconnect layers such as copper layers, aluminum layers, and the like. With the development of semiconductor technology and integrated circuits, semiconductor technology continues to advance to smaller technology nodes, a semiconductor substrate may include an increased number of integrated circuits. The processed semiconductor substrate is then cut into a number of finished integrated circuits that can be individually packaged.\nHowever, with the continuous shrinking of the size of the semiconductor devices, in processing semiconductor devices with technology nodes of 40 nm and below, the finished products after the laser grooving and diamond cutting processes suffer cracking, peeling, and delamination defects.\nTwo methods are currently available to solve delamination problems: the first one is to form an opening in the top layer of the passivation layer in the scribe line region; and the second one is to reduce the size of the test pad and copper metal layer in the scribe line to avoid stress due to laser cutting. However, the conventional methods cannot prevent delamination from occurring.\nFIG. 1 is a cross-sectional view of a semiconductor device containing a sealing ring, as known in the prior art. A semiconductor device (alternatively referred to as integrated circuit or chip) 10 includes semiconductor substrate 100, a chip edge 110, a cutting region disposed at an outer periphery of the semiconductor device, and a seal ring 120 disposed between the cutting region and a chip region to effectively prevent cracking and peeling in the chip caused by dicing of the semiconductor substrate. The seal ring is formed in the metal layer and the dielectric layer of a semiconductor substrate. The seal ring may be a stacked structure of several metal layers connected together through the through-holes (vias) in the interlayer dielectric layers or other mechanical stress buffering structures.\nDelamination occurs because stress is generated when a metal pad is cut during the laser scribing process, the stress is released after transmitting to the chip. Conventional seal rings cannot stop delamination, low dielectric constant materials are used in next-generation semiconductor processes, which greatly reduce the fracture toughness and adhesion between the layers made of different materials, thus expanding the range of the delamination and cracking in the dielectric layers. Delamination and cracking are generally generated at the edge of the chip and extend to the middle of the chip, adversely affecting the functional integrity and yield of the chip.\nThus, there is a need to provide a new semiconductor device that overcomes the disadvantages associated with prior art."}
-{"text": "1. Filed of the Invention\nThe present invention relates to a transfer apparatus providing and a molded article production apparatus having the transfer apparatus. The transfer apparatus provides a component for insert molding to a die by insert molding, and removes a molded article which is formed by the die from the die.\n2. Description of the Related Art\nGenerally, various electronic devices, for example a light such as headlights and a rear light, and a motor such as a starter motor and a motor for air condition, are mounted in a vehicle as a mobile object.\nA junction block is arranged in appropriate place of the vehicle so as to provide electrical power into the above mentioned various electronic devices. The junction block is included by putting together various electronic circuit units of multiple fuses and relays.\nA junction block has a fuse, a relay or a bus bar. Thereby, the junction block is called a fuse block, a relay box, or an electronic connection box. In the present invention, the above fuse block, the relay box and the junction block are called an electronic connection box as a generic name in the below.\nThe above mentioned electric connection box has a box body forming an outer case of it and a wiring block. The box body is made of an insulation synthetic resin, and is formed into a box shape. Electrical part, for example a relay and a fuse, and a connector of a wiring harness are mounted on the box body.\nThe wiring block has a conductive bus bar, which is buried into the box body, and a terminal. The wiring block connects electrically each other the electrical part, which are attached to the box body, and the terminal of the connector of the wiring harness in a predetermined pattern.\nThe bus bar and the terminal constituting the wiring block are set into a cavity of die for insert molding, and the cavity is filled with synthetic resin. Thereby, the bus bar and the terminal are integrally formed with the box body together. That is, the bus bar and the terminal are integrally formed with the box body by insert molding.\nHeretofore, the above mentioned insert molding is formed by various molding production apparatus (for example, see a patent document 1). A molded article production apparatus shown in the patent document 1 has an insert molding machine, a robot as a transfer apparatus providing components into a die and a robot as a transfer apparatus removing a formed molding from the die.    Patent document 1:            Japan published patent application H9-201849        "}
-{"text": "The present invention relates to a process for preparing conveniently an ascorbic acid-2-monophosphate salt, particularly an L-ascorbic acid-2-monophosphate (hereinafter sometimes abbreviated to xe2x80x9cAsMPxe2x80x9d) salt in a high yield, by hydrolyzing an ascorbic acid-2-polyphosphate or a salt thereof, particularly an L-ascorbic acid-2-polyphosphate (hereinafter sometimes abbreviated to xe2x80x9cAsPPxe2x80x9d) or a salt thereof, as the starting material in the presence of magnesium ion.\nL-Ascorbic acid (vitamin C) is known to have various physiological and pharmacological actions, and particularly by virtue of its effect of preventing melamine pigmentation, L-ascorbic acid has been used in whitening cosmetics. However, L-ascorbic acid is unstable to oxygen or heat, therefore, L-ascorbic acid derivatives stable against oxygen or heat have been heretofore formed by converting the hydroxyl group at the 2-position into a phosphate.\nA salt, particularly a magnesium salt of L-ascorbic acid-2-monophosphate (hereinafter the xe2x80x9cmagnesium L-ascorbic acid-2-monophosphatexe2x80x9d is sometimes abbreviated to xe2x80x9cAPMxe2x80x9d) is being used as a stabilized vitamin C derivative. APM exhibits excellent stability in cosmetic materials and scarcely decomposes. Furthermore, APM is easily absorbed through the skin and by the action of phosphatase present inside the human body, L-ascorbic acid is liberated to bring about various physiological actions such as prevention of melamine pigmentation.\nHigh-purity AsMP salts are useful as a stabilized derivative of L-ascorbic acid and can be used in cosmetics, medical products, food additives, feed and other various industrial fields.\nFor monophosphorylating L-ascorbic acid, three methods described below are known. A first method comprises monophosphorylating L-ascorbic acid using phosphorus oxychloride as the phosphorylating agent and this method is described, for example, in JP-B-45-30328 (the term xe2x80x9cJP-Bxe2x80x9d as used herein means an xe2x80x9cexamined Japanese patent publicationxe2x80x9d), JP-B-52-18191, JP-B-59-4438 and JP-A-2-27969 (the term xe2x80x9cJP-Axe2x80x9d as used herein means an xe2x80x9cunexamined published Japanese patent applicationxe2x80x9d) and the like. The second method comprises transposing the phosphoric acid group of a phosphoric acid donor (e.g., adenosine triphosphate, pyrophosphoric acid and the like) to L-ascorbic acid using an enzyme and this method is described, for example, in JP-A-2-283283. The third method comprises reacting a soluble salt of metaphosphoric acid used as a phosphorylating agent with L-ascorbic acid and drying the aqueous solution to obtain AsMP salt and this method is described, for example, in JP-A-5-155893.\nThe first method above is most frequently used and according to Carbohydrate. Res., 67, 127-138 (1978), L-ascorbic acid-2-monophosphate is obtained as a main product in the form of tricyclohexylamine salt in a yield of 86%. The by-products are L-ascorbic acid-3-monophosphate, L-ascorbic acid-2-diphosphate (hereinafter sometimes abbreviated to xe2x80x9cAsDPxe2x80x9d) and 2,2xe2x80x2-bis(L-ascorbic acid)phosphate. For the purification thereof, a complicated step such as ion exchange chromatography is necessary. L-ascorbic acid has four hydroxyl groups active in the reaction with phosphorus oxychloride and phosphorus oxychloride has three active sites, therefore, many by-products are produced. Furthermore, chloride ion generated from phosphorus oxychloride after the reaction is 3 molar times the phosphorus oxychloride used and this requires purification by electrodialysis and the like.\nThe second method is advantageous from the viewpoint that generation of by-products is prevented but still has the defect that the productivity is low. According to JP-A-2-283283, AsMP is produced using an enzyme fixed to an ion exchange resin and using diphosphoric acid (pyrophosphoric acid) as a phosphoric acid donor. However, the yield of AsMP of 34% is very low. Also, the amount of AsMP produced is as low as 1.7% in terms of the concentration in the reaction solution and large scale reaction equipment is necessary. Furthermore, in order to obtain high-purity AsMP salt, purification through many steps such as removal of a large amount of water and recovery of unreacted L-ascorbic acid is required.\nThe third method is a method of drying L-ascorbic acid-2-triphosphate (hereinafter sometimes abbreviated to xe2x80x9cAsTPxe2x80x9d) in an aqueous solution and thereby stepwise hydrolyzing the phosphate bonded to produce AsMP. According to JP-A-5-155893, AsMP is produced by preparing an aqueous solution of AsTP from a soluble salt of metaphosphoric acid and L-ascorbic acid, adjusting the pH thereof to about 5.5 to 6.5 and performing drying at a temperature of from 120 to 180xc2x0 C. As known from the pH of aqueous solution, this stepwise hydrolysis reaction of phosphate is acid hydrolysis reaction. Therefore, although it is set forth therein that sodium and calcium are advantageous as the cation present in the aqueous solution, the alkali property of the cation has no relation to the acid hydrolysis, since the pH of the aqueous solution is acidic.\nFurthermore, the product AsMP obtained by this production method has extremely low purity and the yield is also low. As seen from the composition of the product obtained by this method described in the Examples, the AsMP purity is 61% and the yield based on AsTP is 66%. In addition, much L-ascorbic acid, AsDP and AsTP are present as by-products. The presence of L-ascorbic acid shows that the acid hydrolysis reaction of AsTP extended to L-ascorbic acid over AsMP.\nThis means that it is useless to phosphorylate the hydroxyl group at the 2-position of L-ascorbic acid so as to improve lack of stability of L-ascorbic acid to oxygen or heat. Furthermore, the L-ascorbic acid present in the product decomposes due to oxygen or heat causing coloration or produce oxalic acid which is a skin irritant. The product cannot be used in cosmetics and the like as it is. Therefore, the impurities must be purified by a complicated method such as ion exchange chromatography. For these reasons, this production method of AsMP by the drying of AsTP aqueous solution is not an industrial attractive.\nAs described above, AsMP salt is useful as a stabilized derivative of L-ascorbic acid, nevertheless, the production methods thereof are complicated and disadvantageous in industry. Under these circumstances, a simplified preparation process is demanded therefor.\nAs a result of extensive studies to eliminate the above-described defects, the present inventors have found that high-purity AsMP salt can be conveniently prepared in a high yield by hydrolyzing, preferably under heating AsPP, or a salt thereof at a pH of 7 or more as a starting material in the presence of magnesium ion. The present invention has been accomplished based on this finding of a new process.\nThe present invention provides the following embodiments for preparing an ascorbic acid-2-monophosphate salt:\n(1) a process for preparing an ascorbic acid-2-monophosphate salt, comprising hydrolyzing an ascorbic acid-2-polyphosphate or a salt thereof in the presence of magnesium ion at a pH of 7 or more;\n(2) the process as described in (1) above, comprising hydrolyzing a reaction product of an ascorbic acid and a polyphosphoric acid or a salt thereof in the presence of magnesium ion at a pH of 7 or more;\n(3) the process as described in (1) above, comprising eating a mixture comprising:\n(A) an ascorbic acid-2-polyphosphate or a salt thereof,\n(B) water,\n(C) a magnesium compound capable of supplying magnesium ion in water, and\n(D) a base necessary for maintaining a pH of 7 or more;\n(4) the process as described in (3) above, wherein the\nascorbic acid-2-polyphosphate or a salt thereof (A) is one or more selected from the group consisting of an L-ascorbic acid-2-triphosphate and a salt thereof, an L-ascorbic acid-2-diphosphate and a salt thereof, and a mixture thereof;\n(5) the process as described in (3) above, wherein the magnesium compound (C) is one or more selected from the group consisting of magnesium chloride, magnesium sulfate, magnesium acetate, magnesium nitrate, magnesium oxide, magnesium hydroxide, magnesium carbonate and a mixture thereof;\n(6) the process as described in (3) above, wherein the base (D) is one or more selected from the group consisting of an alkali metal hydroxide, an alkali metal carbonate, an alkaline earth metal hydroxide, an alkaline earth metal oxide, an alkaline earth metal carbonate, a tertiary amine and a mixture thereof;\n(7) the process as described in (3) above, wherein the compound serving both as the magnesium component (C) and the base (D) is one or more selected from the group consisting of magnesium hydroxide, magnesium oxide and magnesium carbonate;\n(8) the process as described in (3) above, wherein the heating temperature is from 35xc2x0 C. to the boiling point of the aqueous reaction solution;\n(9) the process as described in (3) above, which includes, after the completion of reaction, adjusting the pH by adding an acid to the reaction solution or the diluted or concentrated solution to precipitate a phosphate salt and removing the phosphate salt;\n(10) the process as described in (3) above, which includes, after the completion of reaction, adjusting the pH by adding an acid to the reaction solution or the diluted or concentrated solution to precipitate a phosphate salt and dissolving the precipitated phosphate salt and adding a base to the dissolved solution to precipitate the phosphate salt and removing the phosphate salt;\n(11) the process as described in (9) above, which includes crystallizing an ascorbic acid-2-monophosphate salt from the aqueous reaction solution or an aqueous solution after the removal of phosphate using a water-soluble organic solvent; and\n(12) the process as described in (11) above, which includes drying the wet ascorbic acid-2-monophosphate salt crystallized from the aqueous reaction solution or an aqueous solution after the removal of phosphate using a water-soluble organic solvent to obtain the ascorbic acid-2-monophosphate salt as a powder.\nIn the present invention, the ascorbic acid structure may be any of an L-form, a D-form and a racemic compound, however, the present invention is described below by referring to the L-form which is industrially useful.\nAccording to the process of the present invention, an aqueous solution containing AsPP or a salt thereof is hydrolyzed, preferably under heating, at a pH of 7 or more in the presence of magnesium ion and thereby the phosphoric acid ester bond of AsPP is stepwise alkali-hydrolyzed to prepare AsMP.\nThis preparation process of the present invention comprises allowing magnesium ion to be present together and adjusting the pH to 7 or more, the stepwise hydrolysis from AsPP to AsMP proceeds almost constantly in quantity and the hydrolysis reaction from AsMP to L-ascorbic acid scarcely proceeds.\nIn the preparation process of the present invention, the ion allowed to be present at the hydrolysis cannot be just any alkaline earth metal ion but must be magnesium ion. The magnesium ion has an action of accelerating the hydrolysis reaction from AsPP to AsMP but not accelerating the hydrolysis from AsMP to L-ascorbic acid. The reasons therefor are not clearly understood but the magnesium ion maintains relatively mild alkalinity even in a strong alkali condition and has low water-solubility and these properties are considered to take a part in the above-described action.\nFor example, when magnesium oxide was added in 1.7 molar times to an aqueous solution containing AsMP and the solution was heated at 60xc2x0 C. for 6 hours, the amount of residual AsMP was 100 wt %. When the same operation was performed using calcium oxide in place of magnesium oxide, the amount of residual AsMP was 41 wt % and this clearly showed formation of L-ascorbic acid and furthermore formation of oxalic acid by the decomposition of L-ascorbic acid. In the examination where magnesium chloride was added in 1.7 molar times to a solution containing AsMP sodium salt and heated at 60xc2x0 C. for 6 hours, the amount of residual AsMP was also 100 wt %. However, when the same operation was performed using calcium chloride in place of magnesium chloride, the amount of residual AsMP was 63 wt %. The same results were obtained in the examination using barium ion in place of calcium ion.\nAs such, magnesium ion has an action of not accelerating the hydrolysis of AsMP but calcium ion and the like have no such an action. which can be used in the present invention include polyphosphoric acids such as triphosphoric acid, diphosphoric acid (pyrophosphoric acid) and metaphosphoric acid, and salts thereof. Accordingly, the AsPP usually includes AsTP, AsDP, and salts and mixtures thereof. These may be prepared by the method described in U.S. Pat. 5,149,829 and the like. In the present invention, the dehydration product of an ascorbic acid and a polyphosphoric acid or a salt thereof may be used as it is in the hydrolysis.\nAccording to U.S. Pat. 5,149,829, an L-ascorbic acid and a polyphosphate salt such as sodium trimetaphosphate were reacted at a temperature of from 33 to 35xc2x0 C. while adjusting the pH to 10.5 to 10.7 by NaOH or KOH. After the completion of reaction, the inorganic phosphoric acid compound was removed by precipitation as a calcium salt using calcium chloride at a temperature as low as 10xc2x0 C. Thereafter, the residue was concentrated, crystallized from ethanol and dried to obtain the objective AsPP salt as a mixture of AsMP salt, AsDP salt and AsTP salt. In the method of U.S. Pat. 5,149,829, the inorganic phosphoric acid compound precipitated is removed, however, in the present invention, inorganic phosphoric acids may be mixed in the AsPP salt. In the case where use of AsPP as an acid but not as a salt is intended, the acid may be obtained by dissolving an AsPP salt in water and decationizing it by passing the solution through a strongly acidic cation exchange resin at a low temperature. Also, when use of only pure AsTP or only pure AsDP as a raw material is intended, the raw material may be obtained by purifying the AsPP using chromatography.\nIn the present invention, the AsPP concentration is, for example, in water, suitably from about 1 to 80 wt %. If the AsPP concentration is low, the reaction slowly proceeds, whereas it is excessively high, the viscosity increases or magnesium salt of AsPP precipitates. Accordingly, the AsPP concentration is preferably from about 5 to 50 wt %, more preferably from about 10 to 30 wt %.\nThereafter, magnesium ion is added to the aqueous solution of AsPP or a salt thereof. In the case where the AsPP is a salt, a water-soluble magnesium compound such as magnesium chloride is preferably used and in the case where the AsPP is an acid, a magnesium compound showing alkalinity such as magnesium oxide is preferably used. Examples of the magnesium ion-supplying compounds include magnesium chloride, magnesium sulfate, magnesium acetate, magnesium nitrate, magnesium oxide, magnesium hydroxide, magnesium carbonate and the like. These compounds may be used individually or as a mixture without any problem.\nThe amount of magnesium ion added is an equimolar amount or more to the phosphoric acid group of AsPP. As the alkali hydrolysis of AsPP proceeds, phosphoric acid is produced in the aqueous reaction solution and the phosphoric acid produced reacts with magnesium ion present in the aqueous reaction solution to form magnesium hydrogenphosphate or magnesium phosphate. These magnesium salts of phosphoric acid have very low solubility in water and accordingly precipitate out. As a result, magnesium ion in the aqueous solution necessary for the hydrolysis reaction decreases. Accordingly, in order to maintain the magnesium ion in the reaction system even when the hydrolysis reaction of AsPP proceeds, the magnesium ion is required in an equimolar amount or more to the phosphoric acid group of AsPP. In this reaction, the amount of magnesium ion added has no upper limit in particular, however, the presence of excess magnesium ion has no effect on the reaction and this is of no benefit. Accordingly, the amount of magnesium added is preferably from about 1.1 to 2.0 mol, more preferably from 1.3 to 1.8 mol, based on the phosphoric acid group of AsPP.\nIn the preparation process of the present invention, the pH is suitably set to be 7 or more. If the pH is less than 7, the stepwise hydrolysis reaction in the present invention becomes hydrolysis accelerated by an acid, as a result, the hydrolysis of the objective product AsMP is also accelerated to decompose it into L-ascorbic acid and phosphoric acid. Accordingly, the aqueous reaction solution must be set to a pH of 7 or more using an appropriate alkali. Examples of the alkali used here include alkali metal hydroxides, alkali metal carbonates, alkaline earth metal hydroxides, alkaline earth metal oxides, alkaline earth metal carbonates, tertiary amines, and mixtures of two or more thereof. In the case where the alkaline earth metal compound used is magnesium oxide, magnesium hydroxide or magnesium carbonate, magnesium ion can be supplied to the aqueous reaction solution and thus actions of adjusting the pH and at the same time accelerating the hydrolysis of AsPP can be achieved. In the present invention, substantially no problem arises as long as the pH is 7 or more. However, the reaction slowly proceeds in the alkali region closer to the neutral, whereas reaction in the strong alkali condition requires alkali-resistant equipment. Accordingly, the pH is preferably from 8 to 13, more preferably from 9 to 11.\nIn addition, heating is preferably performed in the present invention. If the reaction is performed at a low temperature, a very long time is necessary until the completion of reaction and the plant is not utilized efficiently and this is disadvantageous. Accordingly, the reaction is preferably performed at a temperature of from 35xc2x0 C. to the boiling point of aqueous reaction solution, more preferably from 50 to 90xc2x0 C.\nThe thus-obtained AsMP-containing aqueous solution is diluted or concentrated to an appropriate AsMP concentration and a phosphate salt precipitated (example: phosphoric acid alkali salt) is removed by the solid-liquid separation. At this time, the pH is adjusted with an appropriate acid to effectively precipitate the phosphate salt. Examples of said acid are concentrated or diluted hydrochloric acid. Or, the phosphate salt precipitated may be dissolved by an acid and again precipitated by adding a base.\nIn the aqueous solution containing high-purity AsMP salt after the removal of the phosphate salt precipitated, AsMP salt is precipitated using a water-soluble organic solvent. Examples of the water soluble organic solvent used here include methanol, ethanol, isopropanol, tetrahydrofuran, dioxane, acetone and the like. Among these, ethanol is preferred in view of the cost and ease of recovery.\nThe thus-obtained wet AsMP salt is dried by vacuum drying or fluidized bed drying to obtain AsMP salt in a high yield.\nIn particular, AsMP magnesium salt has excellent stability against heat or light, accordingly, it is used not only for cosmetic powders or lotions but also, for example, for medical products (e.g., oral cavity preparation, eyedrops, products for bath and the like), cosmetics (e.g., lotion, emulsion, cream, pack and the like), food (e.g., bread and the like) and animal feed (e.g., feed for breeding prawn, salmon, yellow tail, eel, carp and the like)."}
-{"text": "The present invention relates to a direct viewing type liquid crystal display device which is employed in a liquid crystal television set, a liquid crystal display for a computer and the like.\nIn recent years, the technical progress of a direct viewing type liquid crystal display device is significant especially in a device employing a color display element. There are many displays having a display quality which is comparable to that of a CRT. In the black-white display, until several years before, the main stream of technology had been a reflecting type liquid crystal display element which did not employ a backlight. However, currently, almost all of them are replaced by a transmitting type liquid crystal display element in use of a backlight even in the black-white display. In the color liquid crystal display, the display can not work without a backlight, and the backlight is an indispensable device in the direct viewing type liquid crystal display device.\nIn a so-called xe2x80x9cnotebook type personal computerxe2x80x9d which has come in use in recent years, the portability is important, and therefore, driving thereof by a battery is indispensable. However, currently, the time capable of driving the device without charging the battery is only several hours, and does not reach a level wherein a day\"\"s operation can continuously be performed. From this viewpoint, the prolongation of the time of continuous use is extremely important. Especially, the illumination device is a device consuming much power in that system, and reducing the power consumption of the illumination device is of great significance.\nIn the mean time, there is a specific contrast ratio distribution in accordance with the viewing angle in a liquid crystal display element employed in the notebook type personal computer. A representative example is shown in FIG. 5 in case of a super twisted nematic liquid crystal display element. FIG. 5 shows the representative viewing cone of the super twisted nematic liquid crystal display element and a diagram designating a change in the contrast ratio of a section on a horizontal line passing through the center of the contrast ratio of the viewing cone. According to FIG. 5, the viewing angle of the liquid crystal display element is widened from a direction perpendicular to the image plane substantially by 40xc2x0 through 50xc2x0, and it is revealed that there is a region having an especially high contrast ratio in the vicinity of the center.\nIn an actual use, the device is often set such that the contrast ratio is maximized on the image plane viewed from an operator in case of the notebook type personal computer. Accordingly, the illuminating efficiency can substantially be promoted, when the maximum brightness of illumination is formed in the viewing direction (mostly a direction perpendicular to an image plane or a direction a little deviated therefrom) which causes the maximum contrast ratio.\nThe viewing angle of the liquid crystal display element employed in the notebook type personal computer is widened substantially to not less than 40xc2x0 through 50xc2x0, and investigations are being performed for widening the viewing angle. Therefore, it is important to adjust the light direction distribution of the illumination device so that the display can be viewed from a more or less oblique direction.\nThe adapting of the brightness distribution of the illumination device to this contrast ratio distribution in such a manner is significant as a means for promoting substantial brightness.\nThe color liquid crystal display device is grossly classified into two systems, namely, a twisted nematic (TN) liquid crystal display device by the active matrix driving using TFTs and a super twisted nematic (STN) liquid crystal display device by the multiplex driving. Both are provided with a construction wherein polarizing sheets are disposed on the light-incident side and the light-emitting side of an element wherein the liquid crystal layer is hold by glass substrates, and the liquid crystal display system is operated by modulating the polarization state of an incident linearly polarized light.\nHowever, the direction of polarization of the incident light of the conventional liquid crystal display element, is not uniform due to the randomly polarized light. Therefore, more than half of the incident light is absorbed by the polarizing sheet disposed on the light-incident side of the display element, and does not substantially contribute to the illuminating light. It is proposed as a structure for reusing the light to be absorbed by the polarizing sheet, wherein a polarized light separator is interposed between a light source lamp and a liquid crystal display device in a projecting type liquid crystal display element, for separating a non-polarized light to mutually orthogonal polarized lights, one polarized light is directly emitted from the polarized light separator, and the other polarized light is converged to the light source lamp and is reused as a light source light (Japanese Unexamined Patent Publication No. 184429/1992).\nHowever, this method is performed on the premise of a projector (projecting type), and it is necessary to provide a sufficient distance between the light source and the polarized light separator. Further, the device effectively functions as an illumination for the projecting type liquid crystal display element, only when the light is a considerably collimated parallel ray. Accordingly, the device is not suitable to adopt as the illumination for a direct viewing type display element wherein thinning thereof is an indispensable condition, and the brightness distribution of the illumination device should adapt to the contrast ratio distribution of the liquid crystal display.\nFurther, it is proposed that a prism array is interposed between a light source for illumination and a display element, as a means for converging light in a direction perpendicular to a display face. However, according to this means, the brightness in the direction perpendicular to the display face is promoted by narrowing the direction of an illuminating light in a specified range, and therefore, the light direction distribution of the illuminating light is narrowed. Further, the brightness in the direction perpendicular to the display face is not sufficient even by this means. Accordingly, an illuminance distribution suitable for the direct viewing type liquid crystal display element can not be provided.\nIn this invention, attention has been paid to the fact wherein only a polarized light having a specific direction of polarization contributes to the promotion of the illuminance of the directly viewing type liquid crystal display element, and a light in a specific direction is selectively converted its polarization direction, among polarized lights which do not contribute to the promotion of the illuminance of the liquid crystal display element by themselves. In this way, the intensity of light can be enhanced in the specific direction while maintaining a wide light direction distribution, with respect to a polarized light having a direction of polarization which can contribute to the promotion of the illuminance, and the light distribution is suitable for the directly viewing type liquid crystal display.\nAccording to a first aspect of the present invention, there is provided an illumination device for a direct viewing type display element comprising:\na flat light guide;\na light source set such that light is incident on a side portion of said flat light guide;\na polarized light separating sheet set on a first side of a light emitting side of the flat light guide for transmitting a p polarized light component and reflecting at least a portion of an s polarized light component with respect to a light ray substantially having a predetermined direction of incidence; and\na light reflecting sheet disposed on a second side opposite to said light emitting side of the flat light guide in parallel with the light emitting site.\nAccording to a second aspect of the present invention, there is provided the illumination device for a direct viewing type display element according to the first aspect, wherein the polarized light separating sheet is comprising a multi-layered structure wherein light transmitting media having a relatively large refractive index and light transmitting media having a relatively small refractive index are laminated.\nAccording to a third aspect of the present invention, there is provided the illumination device for a direct viewing type display element according to the first aspect, wherein the polarized light separating sheet comprises a transparent supporter and at least one dielectric thin film laminated on said transparent supporter having a thickness which is equal to or smaller than a wavelength of visible light.\nAccording to a fourth aspect of the present invention, there is provided the illumination device for a direct viewing type display element according to the first aspect, wherein the polarized light separating sheet comprises a plurality of laminated transparent polymer layers having different refractive indices.\nAccording to a fifth aspect of the present invention, there is provided a liquid crystal display device, wherein the illumination device according to Claim 1 is disposed on a rear side of a direct viewing type liquid crystal display element such that a principle polarization direction of emitted light from the illumination device substantially agrees with a direction of an optical axis of polarization of a polarizing sheet on a light-incident side of a liquid crystal display element.\nAccording to a sixth aspect of the present invention, there is provided the liquid crystal display device according to the fifth aspect further comprising:\na light reflecting means disposed between the polarized light separating sheet and the liquid crystal display element for deflecting a direction of a light ray maximizing a light intensity among light distributing directions to a direction substantially perpendicular to a display face of the liquid crystal display element.\nAccording to a seventh aspect of the present invention, there is provided the liquid crystal display device according to the fifth aspect, further comprising:\na means for rotating polarization direction disposed between the illumination device and the liquid crystal display element for rotating the principle direction of emitted light.\nAccording to an eighth aspect of the present invention, there is provided an illumination device for a direct viewing type display device comprising:\na flat light emitting means for emitting a diffused light including a first polarized light component having a first direction of polarization and a second polarized light component having a second direction of polarization perpendicular to the first direction of polarization; component and\na polarized light converting means disposed in an optical path of a light emitted from said flat light emitting means for emitting said first polarized light component and for emitting at least a portion of said second polarized light component after converting said portion of the second polarized light component into the first polarized light component selectively with respect to a light ray substantially having a light direction maximizing a brightness thereof."}
-{"text": "1. Technical Field\nThe present invention relates to a method of fabricating a copper indium gallium selenide (CIGS) thin film for solar cells through co-vacuum evaporation and a CIGS thin film for solar cells fabricated by the same. More particularly, the present invention relates to a CIGS thin film for solar cells fabricated using a simplified co-vacuum evaporation process without deterioration in physical properties of crystal growth of the thin film and band-gap grading by Ga distribution while reducing process costs through simplification of process steps and adjustment of evaporation elements.\n2. Description of the Related Art\nRecently, importance on development of next-generation clean energy sources has increased due to the depletion of fossil fuel reserves. Thereamong, a solar cell is a device that converts solar energy directly into electricity. Solar cells can serve as an energy source to solve energy problems in the future, since they do not emit pollutants and rely upon the sun.\nA solar cell may be divided into a variety of kinds depending upon materials used in a light-absorption layer, and the most currently available solar cell is a silicon solar cell. However, since silicon prices have been rising due to shortage of high purity silicon, a thin film type solar cell is drawing attention. The thin film type solar cell is fabricated to a thin thickness and thus contributes to less material consumption and light weight, thereby providing a wide application range. Studies have been actively made as to amorphous silicon, CdTe, CIS or CIGS as a material for such a thin film solar cell.\nCIS or CIGS is a I-III-VI compound semiconductor material and has higher conversion efficiency than any other thin-film solar cell materials experimentally prepared. Particularly, since the CIS or CIGS can be fabricated to a thin thickness of 10 microns or less and can be stably operated even after long term use, the CIS or CIGS is expected to be a low-cost, high-efficiency solar cell material as an alternative to silicon.\nParticularly, a CIS thin film is a direct transition type semiconductor that can be made thinner, is suitable for optical-conversion due to a band-gap of 1.04 eV, and exhibits high optical-absorption coefficient. CIGS is an alternative to CIS in which a portion of In is substituted with Ga or Se is substituted with S in order to improve low open circuit voltage of CIS.\nA CIGS thin film can be generally fabricated by vacuum deposition or non-vacuum coating. Particularly, vacuum deposition may include co-evaporation, in-line evaporation, a two-step process (precursor-reaction), and the like. Thereamong, co-evaporation has been traditionally used to fabricate a high-efficiency CIGS thin-film solar cell. However, co-evaporation has a problem of difficulty in commercialization due to a complicated process and difficulty in fabrication of a large area solar cell. To solve this problem, a two-step (deposition/selenization) process facilitating mass production has been developed.\nAmong current CIGS deposition processes, a high efficiency technique having an efficiency of 20% or more is a three-step co-vacuum evaporation process developed by NREL wherein, in a first step, In, Ga, and Se are deposited, in a second step, Cu and Se are deposited excessively relative to a normal amount to promote crystal growth, and in a third step, In, Ga, and Se are further supplied to a thin film to adjust a composition of the CIGS thin-film. Particularly, a Cu\u2014Se compound formed in the second step has a low melting point of about 523\u00b0 C. or less, thereby causing crystal growth of the CIGS thin-film through induction of liquid-phase sintering of the compound. Thus, it has been known that a V-type distribution of Ga in the composition of the thin-film is obtained through such a three-step process, thereby providing efficiency improvement through band-gap grading.\nHowever, considering mass production, since regulation of flux of elements in each step is complex and thin-film deposition requires a long time, the three-step process has a drawback of high process cost. Therefore, there is a need for a method of fabricating a CIGS thin film for solar cells through a simplified process capable of realizing crystal growth and band-gap grading by Ga composition distribution, while simplifying process steps and significantly reducing a film-deposition time."}
-{"text": "Unless otherwise indicated herein, the materials described in this section are not prior art to the claims in this application and are not admitted to be prior art by inclusion in this section.\nElectric motor actuators for robotic and automation systems often require a transmission (speed reducer) in order to operate within the speed-torque requirements of a specific application and of the motor. Commonly used systems include multi-stage gearboxes, timing belts, cables with capstans, harmonic drives, and cycloid gearboxes. These systems may have shortcomings, such as too inefficient, susceptible to overload damage, heavy, and require expensive precision manufacturing. Such systems are also often too expensive for consumer products when high performance is required.\nAs an example, harmonic drive systems can be used in high performance applications where low backlash and gear-ratios greater than 50:1 are required. The harmonic drive is proprietary, heavy, inefficient, susceptible to damage from shock load, and expensive for consumer application. As another example, cable drive systems can be lightweight and efficient; however, achieving useful transmission ratios may lead to complex multi-stage designs that require high preload forces, large cable bend radii, and challenging cable management. Often it is desired to integrate a torque sensor such as a strain gauge load cell at or near the output of the transmission in order to achieve closed loop torque control. Practically, such as a case on a rotating output hub, integration of this sensor can prove challenging as the sensor wires typically rotate with the transmission output, and therefore, require cable management, adding complexity and increasing the number of failure points."}
-{"text": "1. Technical Field\nThe present disclosure relates to a printed circuit board (PCB) assembly, especially to a PCB assembly for automatically detecting a computer system the PCB assembly is coupled to.\n2. Description of Related Art\nA PCB, such as a motherboard usually has multiple functions for coupling to different computer systems. The PCB may have different soft configurations to couple to these different computer systems. For example, a server motherboard may alternatively used in a 1 U (unit) server chassis or in a 2 U server chassis. A server chassis the server motherboard used is often identified by human. The corresponding configuration files may be set by adjusting a jumper on the server motherboard manually. It is time consuming to manually identify and adjust the computer system the PCB coupled in."}
-{"text": "1. Field of the Invention\nThis invention relates to poly(ethylene glycol)-poly(orthoester), poly(ethylene glycol)-poly(orthoester)-poly(ethylene glycol), and poly(orthoester)-poly(ethylene glycol)-poly(orthoester) block copolymers.\n2. Background to the Invention\n1. Micellar System for Tumor Targeting\nOne of the major problems in treating cancer is the difficulty of achieving a sufficient concentration of an anticancer agent in the tumor. This is due to the toxicity, sometimes extreme, of such agents which severely limits the amounts that can be used. However, a major discovery in cancer chemotherapy has been the so-called EPR (enhanced permeation and retention) effect. The EPR effect is based on the observation that tumor vasculature, being newly formed vasculature, has an incompletely formed epithelium and is much more permeable than established older vasculature which is essentially impermeable to large molecules. Further, lymphatic drainage in tumors is very poor thus facilitating retention of anticancer agents delivered to the tumor.\nThe EPR effect can be used in cancer targeting by using delivery systems containing anticancer drugs that are too large to permeate normal vasculature, but which are small enough to permeate tumor vasculature, and two approaches have been developed. In one approach, a water-soluble polymer is used that contains an anticancer drug chemically bound to the polymer via a hydrolytically labile linkage. Such drug-polymer constructs are injected intravenously and accumulate in the tumors, where they are internalized by the cells via endocytosis and released in the lysosomal compartment of the cell via enzymatic cleavage of the labile bond attaching the drug to the polymer. Two disadvantages of this approach are that, first, nondegradable, water-soluble polymers have been used, and this requires a tedious fractionation of the polymer to assure that the molecular weight of the polymer is below the renal excretion threshold, and, second, the drug must be chemically attached to the polymer, which in effect creates a new drug entity with consequent regulatory hurdles that must be overcome. The use of polymer conjugates in cancer diagnosis and treatment is discussed in R. Duncan et al., \"The role of polymer conjugates in the diagnosis and treatment of cancer\", S.T.P. Pharma Sciences, 6(4), 237-263 (1996), and an example of an alginate -bioactive agent conjugate is given in Al-Shamkhani et al., U.S. Pat. No. 5,622,718.\nAn alternate approach has been described. In this approach, an AB or ABA block copolymer is prepared where the B-block is hydrophobic and the A-block is hydrophilic. When such a material is placed in water, it will self-assemble into micelles with a hydrophobic tore and a hydrophilic shell surrounding the core. Such micelles have a diameter of about 100 nm, which is large enough that when they are injected intravenously, the micelles can not leave the normal vasculature, but they are small enough to leave the vasculature within tumors. Further, a 100 nm diameter is too small to be recognized by the reticuloendothelial system, thus enhancing micelle lifetime within the blood stream. Additionally, when the hydrophilic block is poly(ethylene glycol), further enhancement of circulation time is noted, as has been observed with \"stealth\" liposomes. The use of block copolymer micelles is reviewed in G. S. Kwon et al., \"Block copolymer micelles as long-circulating drug delivery vehicles\", Adv. Drug Delivery Rev., 16, 295-309 (1995).\nSakurai et al., U.S. Pat. Nos. 5,412,072 and 5,693,751, and Yokoyama et al., U.S. Pat. Nos. 5,449,513 and 5,510,103, describe block copolymers useful as micellar delivery systems where the hydrophilic block is poly(ethylene glycol) and the hydrophobic blocks are various derivatives of poly(aspartic acid), poly(glutamic acid) and polylysine. U.S. Pat. Nos. 5,412,072 and 5,693,751 describe an approach where drugs have been chemically attached to the hydrophobic segment; while U.S. Pat. Nos. 5,449,513 and 5,510,103 describe an approach where hydrophobic drugs have been physically entrapped within the hydrophobic portion of the micelle. This latter approach is clearly preferable because no chemical modification of the drug is necessary.\n2. Bioerodible Block Copolymer Matrix for Controlled Drug Delivery\nIn AB, ABA, or BAB block copolymers comprising a hydrophilic A block and a hydrophobic B block, the A and B blocks are incompatible and on a microscopic scale will phase-separate. This phase separation imparts unique and useful thermal properties to the material.\nThere is considerable prior art in the development of block copolymers comprised of poly(ethylene glycol) and bioerodible hydrophobic segments such as poly(L-lactic acid), poly(L-lactic-co-glycolic acid) copolymers and poly(.epsilon.-caprolactone), and discussion of their use as drug delivery agents. For example, see W. N. E. Wolthuis et al., \"Synthesis and characterization of poly(ethylene glycol) poly-L-lactide block copolymers\", Third European Symposium on Controlled Drug Delivery, 271-276 (1994), L. Youxin et al., \"Synthesis and properties of biodegradable ABA triblock copolymers . . . \", J. Controlled Release, 27, 247-257 (1993), and Bezwada et al., U.S. Pat. No. 5,133,739.\nPoly(orthoesters) are known as potential vehicles for sustained release drug delivery. See, for example, J. Heller, \"Poly (Ortho Esters)\", Adv. Polymer Sci., 107, 41-92 (1993), and references cited therein, Heller et al., U.S. Pat. Nos. 4,304,767, 4,946,931, and 4,957,998, and PCT International Publication No. WO97/25366.\nThe disclosures of these and other documents referred to in this application are incorporated herein by reference.\nHowever, no block copolymer systems have been described where the hydrophobic, bioerodible segment is a poly(orthoester)."}
-{"text": "1. Field of the Invention\nThe present invention relates to a rotation angle-adjustable rotating device, more particularly to a rotation angle-adjustable rotating device for use in an electrical appliance.\n2. Description of the Related Art\nA conventional rotating device is shown to be employed in an electrically operated fan, and includes coaxially disposed first and second gears 12, 13, a third gear 14 formed with a coupler shaft 15, a connecting plate 16, and a swing arm 17. The first gear 12 is driven by a threaded portion 111 of a drive shaft 11 of the fan, while the second gear 13 is meshed with the third gear 14. One end of the connecting plate 16 is securely mounted on the coupler shaft 15 for co-rotation therewith, and the other end thereof is pivoted to a first end of the swing arm 17. A pushing element 171 is fixed on a second end of the swing arm 17, and extends into an elongated slot 181 formed in a fan housing 18 of the fan in such a manner that rotation of the drive shaft 11 results in synchronous rotation of the coupler shaft 15, which, in turn, results in a back-and-forth swing action of the swing arm 17, thereby driving the fan housing 18 to swing through the pushing element 171.\nOne disadvantage of the conventional rotating device resides in that the rotation angle of the fan housing 18 is unadjustable, thereby inconveniencing the user when he or she desires a larger swinging angle of the fan."}
-{"text": "The present invention broadly relates to compositions useful in separating metal ions or other target material including, but not limited to, polypeptides, nucleic acids, and endotoxins, from non-target material."}
-{"text": "Example embodiments of the present inventive concept relate generally to integrated circuit devices and, more particularly, to embedded refresh controllers and memory devices including the embedded refresh controllers.\nSemiconductor memory devices for storing data may be classified into volatile memory devices and non-volatile memory devices. Volatile memory devices, such as dynamic random access memory (DRAM) devices, may be configured to store data by charging or discharging capacitors in memory cells, and lose the stored data when power is off. Non-volatile memory devices, such as flash memory devices, may maintain stored data even when power is off. Volatile memory devices are widely used as main memories of various apparatuses, and non-volatile memory devices are widely used for storing program code and/or data in various electronic devices, such as computers, mobile devices, etc.\nIn volatile memory devices, cell charges stored in a memory cell may be lost by a leakage current. In addition, when a wordline is transitioned frequently between an active state and a precharged state (i.e., when the wordline has been accessed intensively or frequently), an affected memory cell connected to a wordline that is adjacent to the wordline frequently accessed may lose stored charges. Charges stored in a memory cell may be recharged before data is lost by leakage of cell charges. Such recharge of cell charges is referred to as a refresh operation, and a refresh operation may be performed repeatedly before cell charges are significantly lost."}
-{"text": "The invention relates to a lower portion of an electric and electronic enclosure, comprising a bottom and fastening arrangements for fastening a cover to the lower portion with a bayonet-like fastening, and the fastening arrangements comprise a bore and a shoulder-like stopper for receiving a mounting bracket of a spindle-like element.\nThe invention also relates to an electric and electronic enclosure having a lower portion with a bottom, a cover, and fastening means for fastening the cover to the lower portion with a bayonet-like fastening, the fastening means comprising a bore and a shoulder-like stopper arranged to the lower portion, a through-hole, and a stop face formed on the cover, and a spindle-like element that comprises a top end having a shoulder for resting against the stop face of the cover, and at a distance from the top end a mounting bracket for fastening at the shoulder-like stopper of the lower portion to provide a bayonet-like fastening.\nElectric and electronic enclosures of the mentioned type and their lower portions are well known. An advantage of a bayonet-like fastening is that the cover of the enclosure can be quickly and easily fastened to the lower portion of the enclosure. A bayonet-like fastening also has the advantage that it can easily be accomplished in such a manner that the cover can flexibly give way when a possible electric arc ignites inside the enclosure, thus preventing any great damage. One drawback with the bayonet-like fastening is that it is not suitable for situations in which the cover of the enclosure needs to remain fixed to the lower portion. In such a case, ordinary screws are usually used to fasten the cover to the lower portion. In addition, the bayonet screw of the bayonet-like fastening has a complex structure, which makes it difficult and expensive to manufacture in comparison with a conventional screw. A further drawback with a bayonet screw is that it cannot be machine-mounted on the enclosure.\nDepending on whether the enclosure needs to be such that the cover is fastened to the lower portion with a bayonet-like fastening or with ordinary screws, the user has to select the lower portion of the enclosure to suit the end use. Thus, there must be two different lower portions depending on the manner in which the cover is to be fastened: a first type of the lower portion comprises bores with a shoulder-like stopper for receiving the mounting bracket of the bayonet screw, whereas a second type of lower portion comprises a bore with threads for receiving the threads of an ordinary screw. The need to manufacture and store lower portions of two different types is naturally a disadvantageous solution that correspondingly increases costs. Sometimes, it is necessary to change the fastening method of the enclosure cover (from a bayonet-like fastening to a screw fastening or vice versa) as the application of the enclosure changes."}
-{"text": "Services that provide users with guidance information on public transportation services in the form of automated sounds are in wide use. For example, announcement systems used in buses provide passengers with guidance information by way of voice guidance on, for example, bus stops, bus fares, and the like, which guidance voice is played at an appropriate timing, for example at each of bus stops, upon operation of an operator, with such voice guidance being pre-prepared for each of bus routes. Patent Document 1 discloses a configuration in which information such as names of bus stops is announced to passengers through the generation of voice signals that correspond to voice data based on operations such as an input operation into a device installed in a bus carriage. Patent Document 2 discloses a configuration for the creation of voice data corresponding to a guidance voice whereby, supposing that a content of voice guidance is a sentence such as \u201cthe next stop is [name of a bus stop]\u201d, names of bus stops are created by utilizing low-compression coding, while common words such as \u201cnext\u201d, \u201cis\u201d, and the like are created by utilizing high-compression coding."}
-{"text": "1. Field of the Invention\nThe present invention relates to an airbag apparatus which is installed on a vehicle and is adapted to be deployed and inflated between the vehicle and an occupant thereof for protection of the occupant when an impact of a predetermined value or greater is applied to the vehicle when the vehicle is involved in a collision.\n2. Related Art\nConventionally, as an airbag for the driver's seat, there is proposed an airbag apparatus described in Patent Publication No. WO 2005/044643A, for example. The airbag apparatus described in WO 2005/044643A includes an inflator functioning as a gas generator and an airbag and is designed to be deployed and inflated momentarily by the pressure of a gas injected from the inflator when an impact of a predetermined value or greater is applied to a front end of the vehicle when the vehicle is involved in a collision. Namely, in this airbag apparatus, the airbag is brought into abutment with an annular rim portion of a steering wheel so as to be deployed and inflated between the steering wheel and an occupant (the driver) seated in the driver's seat when the airbag is deployed and inflated. In addition, when the occupant is captured by the airbag to press the bag toward the front of the vehicle due to the vehicle being involved in a collision, the airbag is not only stopped and received but also is given a reaction force directed to the rear of the vehicle by the rim portion, so that the occupant is protected by the airbag.\nIncidentally, in recent years, steer-by-wire type steering systems (steering wheels) which do not use the column shafts draw attentions. In the steering systems like this, there occurs a case where in place of a normal steering wheel in which a rim portion is formed into an annular shape, an irregular steering wheel is used in which part of a normally annular shaped rim portion is open (hereinafter, referred to as a \u201cirregular steering wheel\u201d). In this irregular steering wheel, when looked about in a radial direction from a pad portion where an airbag apparatus is installed, it appears that part of the rim portion is open. In addition, as the irregular steering wheels, an irregular steering wheel is known in which a rim portion is made up of a pair of left and right grip portions which are each formed into a minor arc and the overall shape of the steering wheel including both the grip portions and a connecting portion which connects both the grip portions together is formed into substantially an H-shape as viewed from the top.\nIn the event of the irregular steering wheel like this, however, when an airbag is deployed and inflated between both the grip portions (namely, at an opening portion of the rim portion), the airbag needs to obtain a reaction force (a reaction force directed to the rear of a vehicle) which is given thereto when restraining (protecting) the occupant by any other means than the rim portion."}
-{"text": "Bone grafts are commonly used to treat defects in the skeletal system caused by injury or disease. Skeletal defects often require bone grafts to maintain space and provide a matrix for healing. A graft should provide or facilitate the various mechanisms of bone healing including osteoconduction, osteoinduction, and osteogenesis. Osteoconduction is the ability of the graft to act as a matrix or scaffold to support bone formation. Osteoinduction is a result of bone growth factors that stimulate differentiation of local cells to become bone forming cells, i.e. osteoblasts. Bone morphogenic proteins (BMP's) that are naturally occurring in bone, or that may be produced by recombinant gene technologies, are responsible for osteoinduction. Osteogenesis refers to the formation of bone, and may also be used to reference the ability of cells, to form bone. Bone forming cells may either be resident at the graft site or transplanted to the site by autogenous bone, bone marrow aspirate and/or cell implantation. Considering these requirements to form bone, a need exists for a reproducible and cost-effective process of making a bone graft having improved osteoconductive and osteoinductive properties."}
-{"text": "Additionally, in response to increasing social demands for energy and resource saving, there is an increasing need for a rubber material that is excellent in weather resistance and crack growth resistance to meet the requirements for improving durability of tires in order to reduce fuel consumption of automobiles. In addition, in view of the recent surge in price of butadiene, it is expected that the price of raw materials will even more dramatically rise in the future. As such, there is a need to make use of inexpensive olefin resource even in tire materials. To address this issue, some techniques have been proposed to improve weather resistance by blending an ethylene-propylene-diene rubber (EPDM) containing olefin as a principal component. These techniques, however, suffer from drawbacks, such as lower fracture resistance and lower crack growth resistance (see, for example, JP 2000-063639 A (PTL 1)).\nIn addition, applying a copolymer of a conjugated diene and a non-conjugated olefin to the compounded rubber results in fewer double bonds in the conjugated diene unit (conjugated diene compound-derived unit) in the copolymer as compared with conjugated polymers, and therefore, improved ozone resistance. In addition, the characteristics other than ozone resistance required when a rubber composition is applied to various applications (such as tires, conveyor belts or anti-vibration rubber) include good fracture resistance and good crack growth resistance. It is well known that coordination anionic polymerization using catalyst systems represented by a Ziegler-Natta catalyst allows for homopolymerization of olefins and dienes. However, it was difficult to provide efficient copolymerization of olefins and dienes using such polymerization systems.\nTo address this issue, for example, JP 2000-154210 A (PTL 2) discloses such a catalyst for polymerizing a conjugated diene that contains a transition metal compound of group IV of the periodic table having a cyclopentadiene ring structure, and also discloses \u03b1-olefins such as ethylene as monomers which can be copolymerized with the conjugated diene. However, PTL 2 does not provide a specific description of copolymerization of a conjugated diene compound and a non-conjugated olefin. Obviously, there is no description or suggestion that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nFor example, JP 2006-249442 A (PTL 3) discloses a catalyst for polymerizing olefins that consists of a transition metal compound such as a titanium compound and a co-catalyst, and also discloses a copolymer of an \u03b1-olefin and a conjugated diene compound. However, there is no description or suggestion in PTL 3 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nIn addition, JP 2006-503141 A (PTL 4) discloses a copolymer of ethylene and butadiene that is obtained by synthesizing ethylene and butadiene as a starting material by means of a special organic metal complex as a catalytic component. However, there is no description or suggestion in PTL 4 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nIn addition, JP 2000-086857 A (PTL 5) discloses a butadiene polymer having cis content of 92%, vinyl content of 6% and ethylene content of 3% or 9%. However, there is no description or suggestion in PTL 5 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nIn addition, JP 2000-154279 (PTL 6) discloses a rubber composition including: a butadiene-ethylene block copolymer having a cis content of 92% and ethylene segments in an amount of 4.8 mass % of the total; polybutadiene having a cis content of 95.2% and a vinyl content of 2.5%; and carbon black N220. However, there is no description or suggestion in PTL 6 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nIn addition, JP 11-228743 (PTL 7) discloses an olefin-rich, olefin-diene copolymer. However, there is no description or suggestion in PTL 7 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber."}
-{"text": "1. Technical Field\nThe present invention relates generally to a cloth-like heating element with two pairs of electrical conductors and parallel circuits which are formed via a severing process.\n2. Related Art\nAn existing cloth-like heating element, Korean Utility Model Registration #1100067, is issued to the same applicant as the present invention, and its functions are as follows:\nAccording to predetermined design, a great number of resistance heating wires are arranged at some intervals between individual threads which wind up on the loom's beam, while a fine copper wire and a thread are used as the weft.\nThrough the weaving, the resistance heating wires and the copper wires are interlaced with each other to form the heating element, thereby the cloth-like heating element with parallel structure is completed.\nThe above mentioned cloth-like heating element, which has the benefits induced by parallel circuitry, is only in use with very low voltages (Dc12V, 24V), because in case of increasing the number of heating wires for the purposes of enhancing durability or heat balance of the heating element, the input voltage for the heating element must be relatively decreased.\nThis is in accordance with Ohm's law, which states that electrical resistance is inversely proportional to the number of heating wires connected in parallel.\nTherefore, the above mentioned heating element cannot function under high voltages (AC100V, 220V), which is a fatal weakness."}
-{"text": "Substrates that include one or more layers of semiconductor material are used to form a wide variety of semiconductor structures and devices including, for example, integrated circuits (ICs) (e.g., logic processors and memory devices), radiation emitting devices (e.g., light emitting diodes (LEDs), resonant cavity light emitting diodes (RCLEDs), and vertical cavity surface emitting lasers (VCSELs)), radiation absorbing devices (e.g., optical sensors and solar cells) and switching/rectifying devices (e.g., power electronic devices). Such devices may comprise bipolar junction transistors (BJTs), power metal-oxide-field-effect-transistor (MOSFETS), thyristors, schottky diodes, junction field effect transistors (JFETs), insulated gate bipolar transistors (IGBTs) and PIN diodes. Such semiconductor devices are conventionally formed in a layer-by-layer manner (i.e., lithographically) on a semiconductor substrate.\nHistorically, a majority of such semiconductor substrates that have been used in the semiconductor device manufacturing industry have comprised thin discs or \u201cwafers\u201d of silicon material. Such wafers of silicon material are fabricated by first forming a large generally cylindrical silicon single crystal ingot, and subsequently slicing the single crystal ingot perpendicularly to its longitudinal axis to form a plurality of silicon wafers. Such silicon wafers may have diameters as large as about thirty centimeters (30 cm) or more (about twelve inches (12 in.) or more). Although silicon wafers generally have thicknesses of several hundred microns (e.g., about 700 microns) or more, only a very thin layer (e.g., less than about three hundred nanometers (300 nm)) of the semiconductor material on a major surface of the silicon wafer is actually used to form active devices on the silicon wafer.\nIt has been discovered that the speed and power efficiency of semiconductor devices can be improved by electrically insulating the portion of the semiconductor material that is actually used to form the semiconductor devices from the remaining bulk semiconductor material of the substrate. In addition, it has more generally been discovered that providing the semiconductor material used to form the semiconductor devices upon a base substrate, comprising one or more materials, allows for \u201cengineering\u201d of the properties of the semiconductor material.\nAs a result, so-called \u201cengineered substrates\u201d have been developed that may include a relatively thin layer of semiconductor material (e.g., a layer having a thickness of less than about three hundred microns (300 \u03bcm)) disposed upon one or more other material, such as, for example, dielectric material (e.g., silicon dioxide (SiO2), silicon nitride (Si3N4), silicon carbide (SiC), silicon (Si) or aluminum oxide (Al2O3)). Optionally, the layer of dielectric material may be relatively thin (e.g., too thin to enable handling by conventional semiconductor device manufacturing equipment), and the layer or layers of material over which the semiconductor material is disposed (i.e., the base substrate) may be thick enough to enable handling of the engineered substrate by manufacturing equipment.\nA wide variety of engineered substrates are known in the art and may include semiconductor materials such as, for example, silicon (Si), germanium (Ge), silicon carbide (SiC), III-V type semiconductor materials, and II-VI type semiconductor materials.\nFor example, an engineered substrate may include an epitaxial layer of III-V type semiconductor material on a surface of a base substrate, such as, for example aluminum oxide (Al2O3) (which is often referred to as \u201csapphire\u201d). Using such an engineered substrate, additional layers of material may be foamed and processed (e.g., patterned) over the epitaxial layer of III-V type semiconductor material to form one or more semiconductor devices on the engineered substrate.\nWhen a layer of semiconductor material is epitaxially grown at elevated temperatures upon an engineered substrate, lattice strain may be induced in the crystal lattice of the layer of semiconductor material. The strain in the semiconductor material may result from a lattice parameter mismatch between the crystal lattice of the underlying materials of the engineered substrate and that of the semiconductor material being formed thereon (e.g., the underlying material has one or more lattice constants that differ from that of the semiconductor material).\nIn addition, a lattice strain may also be induced in the crystal lattice of the semiconductor material during epitaxial growth at elevated temperature due to a difference in the coefficients of thermal expansion (CTE) exhibited by the respective adjacent materials. For example, if the underlying engineered substrate has a mean CTE which is greater than that of the semiconductor material being grown upon the engineered substrate, the semiconductor material may be grown in a state of tensile strain. Such a state of tensile strain may increase during the growth of the semiconductor material as the layer thickness increases, and may eventually result in the formation of defects in the semiconductor layer. Such defects may include, for example, dislocations and cracks."}
-{"text": "Iontophoresis is a method for introducing ionic substances into a body. The method utilizes direct electrical current to drive the ionized substances, such as drugs, through the intact skin or other body surface. This has proven to be very useful in numerous medical applications. U.S. Pat. Nos. 3,991,755 issued to Jack A. Vernon, et al and 4,141,359 issued to Stephen C. Jacobsen, et al disclose examples of iontophoretic devices and some applications of the devices. The iontophoretic process has been found to be useful in the administration of lidocaine hydrochloride, hydrocortisone derivatives, acetic acid, fluoride, penicillin, dexamethasone sodium phosphate and many other drugs.\nIn iontophoretic devices two electrodes are used. One electrode, called the active electrode, is the electrode at which the ionic substance is driven into the body. The other electrode, called the indifferent or ground electrode, serves to close the electrical circuit through the body. It will be appreciated by those skilled in the art that the active electrode must hold, contain or otherwise have available to it a source of the ionic substance. Thus, in the prior art the active electrode is generally relatively complex compared to the indifferent electrode.\nGenerally, prior iontophoretic drug delivery systems provide a single drug delivery rate. Such rate is obtained by applying a constant iontophoretic current designed to achieve a certain steady-state therapeutic concentration of drug in the body. With the use of such systems, there is a certain delay between the time that the iontophoretic maintenance current is initiated and when the desired therapeutic level of concentration is reached. Such delay may be, for example, thirty minutes from the time the iontophoretic current is initiated. In many cases, however, it is desirable or necessary that the iontophoretic drug reach therapeutic levels relatively fast. For example, where iontophoresis is used to deliver a narcotic pain killer, the patient often cannot tolerate a delay of even fifteen minutes. If the iontophoretic drive current is initially set at a relatively high level in order to encourage the rapid migration of iontophoretic drug into the bloodstream, the system will ultimately reach a steady-state level higher than desired or therapeutically safe. As a result, there is a need for an iontophoretic delivery system wherein therapeutic levels of drug concentration in the blood can be rapidly obtained while at the same time achieving a desirable steady-state maintenance level of administration."}
-{"text": "The present invention relates to a plant for transmitting electric power between a direct voltage network for High Voltage Direct Current (HVDC), and at least two alternating voltage networks connected thereto through a station with each of the stations being adapted to perform transmission of electric power between the direct voltage network and a respective one of the alternating voltage networks.\nIn such plants, according to the prior art, primarily uninsulated overhead lines have been used to implement a direct voltage network interconnecting stations which may be located very large distances apart, for example, 1000 km. However, such overhead lines have a significant disturbing influence upon the country and the living things where they exist. The alternative has been to use cables having an inner conductor surrounded by a thick insulating layer formed by a paper impregnated by oil, but such cables are so costly that they are not a realistic alternative to overhead lines. Experiments have been carried out with cables having a conductor and an insulating layer on a polymer base surrounding the conductor for transmitting High Voltage Direct Current (HVDC), but these experiments have revealed great problems in using this type of cable for transmitting High Voltage Direct Current (HVDC), since space charges are created in the insulating layer which may cause breakthroughs and breakdowns in the cables. The space charges are caused by the changes of polarity which the conductor is subjected to when the feeding direction of the power transmitted between the direct voltage network and an alternating voltage network is changed. A result of these problems is that overhead lines have been used for transmitting High Voltage Direct Current (HVDC) over land, and cables having tailor made insulating sheets have been used as sea cables for transmitting between land areas separated by water."}
-{"text": "Typically, a loudspeaker unit exhibits a low frequency fall-off in its frequency transfer function. In audio systems, it is known to boost low frequencies by having the option of applying a fixed dedicated \u2018extra gain\u2019 at low frequencies (\u2018Loudness\u2019 button). WO9318626A1 describes dynamic adjustment of low frequency gain, where the applied gain is dependent on the content (magnitude) of signal components at low frequencies in the current signal."}
-{"text": "1. Field of the Invention\nThe present invention relates to keyboard operated musical instruments. More particularly, the present invention relates to interfacing an acoustic or electrically amplified piano to a digital computer, storage device, or music synthesizer.\n2. Description of the Prior Art\nSince before the turn of the century, efforts have been made to capture information from acoustic keyboard instruments as the instruments are played. Early devices of this type were predominantly player pianos wherein a performance could be captured in the form of holes on paper that controlled a pneumatic piano control system. The point of such devices was to make one recording on a paper roll of the performance. Thereafter, the performance could be replayed as desired on a player piano by use of a copy of the original paper roll.\nOne more recent approach to capturing keyboard information during a performance employs a switch matrix including a single matrix row coupled to a bus extending underneath a piano keyboard and having a number of matrix columns corresponding to the number of keys on the instrument. The switch matrix includes a contact on each key operable to contact the matrix bus when the key is operated and thus indicate a switch closure. Such device is shown in U.S. Pat. No. 3,604,299, issued to Englund, Sept. 14, 1971.\nThe installation of such device is time consuming and, therefore, expensive. Once installed, it is likely that the action of the piano may be affected. Additionally, the open contact and bus system is subject to performance degradation and lack of reliability as the result of oxidation and contact wear. A similar arrangement is provided in U.S. Pat. No. 4,023,456, issued to Groeschel on May 17, 1977.\nU.S. Pat. No. 4,104,949, issued to Clark on Aug. 8, 1978, provides a series of keys located on the piano above the keyboard and including sensors responsive to keyboard operation. The location of such device above the keyboard eliminates some of the playing area on the keyboard available to the musician and, thus, interferes with normal keyboard operation.\nThe most recent approach to capturing keyboard information in a mechanical or electrical manner is shown in U.S. Pat. No. 4,307,648, issued to Stahnke on Dec. 29, 1971 and in U.S. Pat. No. 4,351,221, issued to Starnes et al. on Sept. 28, 1982. Both such devices use an optical switch including a light source and photo cell. Key operation results in movement of a key or hammer within the piano mechanism that produces a corresponding movement of a flag attached thereto. As a result, a beam of light between the light source and photo cell is broken by interposition of the flag, and information is produced corresponding to key operation.\nSuch devices are installed to the piano action during a time consuming and, therefore, expensive installation procedure. Because there is continual movement of the flag, there is a likelihood that the flag will become dislodged or slip out of alignment and not operate the optical switch properly. Thus, such systems generally require considerable maintenance. Furthermore, the addition of a flag to the piano action might, in some circumstances, affect the action of the piano and, thus, alter or degrade the feel and playing quality of the piano.\nAn additional problem in prior art keyboard interfaces is that of cabling. A standard piano keyboard has eighty-eight keys. Known keyboard interfaces, therefore, require a cable assembly including eighty-nine wires (eighty-eight key wires and one return or ground wire). Although there have been some attempts to reduce this number of wires through various decoding and latching schemes, generally, significant problems are introduced by the prior art devices when routing cables to transfer electronic signals in response to keyboard operation.\nElectronic musical instruments include the well-known electronic music synthesizer, a device that is capable of imitating many known electronic and acoustic instruments, as well as producing its own unique sounds. Electronic music synthesizers are operated by keyboards, or slide controllers or both. One disadvantage of an electronic music synthesizer is that the keyboard generally provided is of a size (that is, number of keys), and has an action or feel unlike that of acoustic pianos. Many musicians have studied on acoustic pianos. A musician who has become accustomed to the feel of a particular keyboard tends to prefer that particular keyboard over other keyboards having a different feel.\nAnother disadvantage of synthesizers is that they introduce an additional keyboard into a playing configuration. For example, during a performance, a musician who plays both an acoustic keyboard instrument and an electronic music synthesizer has to physically shift between the two. Such additional keyboards take up space in the club or theater where the performance is to be held and provide extra equipment to be transported from location to location in accordance with the musician's itinerary."}
-{"text": "Cathodes can emit electrons by photoemission, thermionic emission, and field emission, or as the result of negative electron affinity. A field-emission cathode (or field emitter) provides electrons when subjected to an electric field of sufficient strength. The electric field is created by applying a suitable voltage between the cathode and an electrode, typically referred to as the anode or gate electrode, situated a short distance away from the cathode.\nVarious techniques have been explored for creating field emitters. Chason, U.S. Pat. No. 5,019,003, fabricates a field emitter by depositing preformed electron-emissive objects on a substrate consisting of dielectric and/or electrically conductive material. The preformed objects, which have sharp edges, can consist entirely of electron-emissive material such as molybdenum or titanium carbide. Alternatively, the preformed objects can consist of electrically insulating cores with thin electron-emissive coatings over the insulating cores. The longest dimension of the objects is approximately 1 .mu.m. A bonding layer is employed to bond the objects to the substrate.\nJaskie et al (\"Jaskie I\"), U.S. Pat. No. 5,141,460, discloses a technique in which diamond is used in fabricating a field emitter. Kane et al (\"Kane I\"), U.S. Pat. No. 5,129,850, discloses a related technique for manufacturing a field emitter that utilizes diamond. The fabrication techniques in Jaskie I and Kane I generally entail implanting carbon into a substrate to create diamond nucleation sites and then growing diamond crystallites at the diamond nucleation sites. The resulting regions of diamond crystallites appear to be electron emissive.\nUse of diamond to provide electrons is desirable for a number of reasons. Depending on how it is produced, diamond can have a low work function. This is advantageous because the electric field needed to emit electrons decreases as the work function decreases. Diamond has a low chemical reactivity. In particular, the gases typically present in a sealed vacuum device such as a CRT have little effect on diamond. Also, changes in temperature affect diamond less than most materials used as electron emitters.\nIn Jaskie I and Kane I, the diamond crystallites are grown by chemical vapor deposition (\"CVD\"). While CVD is economically suitable for depositing many materials, diamond CVD is costly because the diamond CVD growth rate is low and a high CVD temperature is needed. The diamond CVD in Jaskie I and Kane I appears too expensive for low-cost volume production of CRTs in flat-panel televisions.\nJaskie et al (\"Jaskie II\"), U.S. Pat. No. 5,278,475, produces a gated field emitter that utilizes diamond crystallites as electron sources. The diamond crystallites are deposited across the upper surface of a supporting structure consisting of a substrate or a patterned layer of conductive/semiconductive material formed on an electrically insulating substrate. A dielectric layer is deposited over the diamond crystallites. A gate (or control) electrode layer, likewise consisting of conductive/semiconductive material, is deposited on the dielectric layer. Openings are formed through the gate electrode and dielectric layer to expose diamond crystallites at selected areas of the supporting structure.\nKane et al (\"Kane II\"), U.S. Pat. No. 5,252,833, discloses a similar gated field emitter in which diamond crystallites provide electrons. The diamond crystallites in Kane II are situated on conductive/semiconductive paths at the bottoms of openings through a dielectric layer and an overlying gate electrode. The diamond crystallites consist of polycrystalline diamond. Taking note of the fact that the (positive) affinity of a material to retain electrons increases the surface work function and thus increases the electric field needed for an electron to escape the material, Kane II indicates that polycrystalline diamond with a (111) crystallographic orientation is particularly useful as an electron source because (111) polycrystalline diamond has a negative electron affinity.\nElectron affinity is an important consideration in choosing an electron source. However, maintaining a negative electron affinity during volume field-emitter production requires special steps. Also, it is not clear that the diamond crystallites in Jaskie II and Kane II will be securely fixed to the underlying material in a manner that permits a control voltage to be suitably impressed on the diamond crystallites. As a result, the gated field emitters of Jaskie II and Kane II may not perform well. It would be advantageous to have an electron-emitting device in which diamond or a related carbon-containing material can be utilized as an electron source and which can be fabricated in a manner that avoids the above-mentioned disadvantages of the prior art."}
-{"text": "The subject matter disclosed herein relates to fuel injectors, and, more particularly, to fuel injectors for gasifiers.\nAn integrated gasification combined cycle (IGCC) power plant includes a gasifier with one or more fuel injectors. The fuel injectors supply a fuel, such as an organic feedstock, into the gasifier along with oxygen and steam to generate a syngas. However, existing fuel injectors may have several passages for the fuel, oxygen, and other materials injected into the combustion chamber, thereby increasing the complexity and cost of the fuel injectors. In addition, existing fuel injectors may have narrow operating windows and thus, may be replaced when operating conditions for the gasification system change."}
-{"text": "The present invention relates to nanopore/nanochannel devices, and more specifically, to reducing the entropy of molecules via agarose gel in nanopore/nanochannel devices.\nNanopore sequencing is a method for determining the order in which nucleotides occur on a strand of deoxyribonucleic acid (DNA). A nanopore (also referred to a pore, nanochannel, hole, etc.) can be a small hole in the order of several nanometers in internal diameter. The theory behind nanopore sequencing is about what occurs when the nanopore is submerged in a conducting fluid and an electric potential (voltage) is applied across the nanopore. Under these conditions, a slight electric current due to conduction of ions through the nanopore can be measured, and the amount of current is very sensitive to the size and shape of the nanopore. If single bases or strands of DNA pass (or part of the DNA molecule passes) through the nanopore, this can create a change in the magnitude of the current through the nanopore. Other electrical or optical sensors can also be positioned around the nanopore so that DNA bases can be differentiated while the DNA passes through the nanopore.\nThe DNA can be driven through the nanopore by using various methods, so that the DNA might eventually pass through the nanopore. The scale of the nanopore can have the effect that the DNA may be forced through the hole as a long string, one base at a time, like thread through the eye of a needle. Recently, there has been growing interest in applying nanopores as sensors for rapid analysis of biomolecules such as deoxyribonucleic acid (DNA), ribonucleic acid (RNA), protein, etc. Special emphasis has been given to applications of nanopores for DNA sequencing, as this technology holds the promise to reduce the cost of sequencing below $1000/human genome."}
-{"text": "This invention relates to area navigation systems, and more particularly to an area navigation system having improved accuracy in the presentation of course deviation information. The invention is described in the particular context of the area navigation systems described and claimed in Assignee's Abnett et al U.S. Pat. No. 3,796,867 entitled, \"ANGLE MODE AREA NAVIGATION COMPUTER\", and in Assignee's Bean U.S. Pat. No. 3,750,942 entitled, \"DISTANCE MODE AREA NAVIGATION COMPUTER\", the disclosures of which patents are incorporated herein to the extent pertinent.\nAn extended description of the concepts of area navigation and its many advantages is set forth in the Abnett et al and Bean patents, and the reader is referred to these patents for more information. Briefly, however, there exists a network of navigation aids called VORTAC's which provide ratio information from which range and bearing information may be derived. In area navigation, the VORTAC range and bearing data are used in conjunction with an arbitrarily selected range and bearing from the VORTAC to a so-called waypoint. From these, an on-board computer generates range and bearing data from the aircraft to the waypoint. Because the location of a waypoint is arbitrary within the range of a particular VORTAC, the pilot is no longer constrained to fly along or near the established airways (defined with respect to a series of VORTAC stations). Instead, one may fly along any convenient course simply by appropriate choice of waypoints. In broad terms, area navigation systems such as those of the Bean and Abnett et al patents facilitate the above by solving a vector triangle defined by the aircraft, the VORTAC, and the waypoint, given the selected waypoint address vector and the measured aircraft to VORTAC vector.\nIn the Abnett et al patent, there is shown an area navigation system in which the aircraft to waypoint vector is used to generate a display of range from the aircraft to the waypoint, and a track angle deviation display, the latter is based on the difference between the actual aircraft to waypoint bearing (the angle component of the aircraft to waypoint vector) and the bearing at which it is desired that the aircraft intercepts the waypoint, i.e. the selected track angle. The angle display is advantageously accomplished by means of a \"left-right\" needle included in the course deviation indicator (or CDI) of the standard VORTAC navigation systems.\nIn the Bean patent, track deviation is displayed on a left-right meter as distance from the selected track.\nEach of the aforementioned types of systems has particular advantages, as described in the respective patents. In both instances, however, one of the most important factors in assuring accurate and useful system operation (apart from good system design) is overcoming a track deviation display instability due to radio propagation defects associated with the VORTAC station, and corresponding errors inherent in the mathematical formulation of the data, and the display mode required for the track error."}
-{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nAn associated service between an Email box and a mobile communication terminal is to bind a user's Email-box number with the user's mobile communication terminal number, so as to make an Email-box user utilize a mobile communication terminal to receive a new-Email reminding message from a bound Email-box, meanwhile the mobile communication terminal may also be utilized to complete various kinds of Email-box operations, such as reading operation, writing operation, replying operation and forwarding operation about an Email.\nFIG. 1 is a schematic diagram illustrating existed exemplary interactive operations between a mobile communication terminal and an Email box. With reference to FIG. 1, the core of interactive operations between the mobile communication terminal and the Email box lies in transit server. The transit server generally needs to complete the following functions.\n1) Receiving a new-Email reminder sent by the Email-box server via the Internet, filtering according to filter-conditions set by user, sending a new-Email reminder which meets the filter-conditions to the mobile communication terminal through Short Message Service (SMS).\n2) The mobile communication terminal initiates requests, such as reading-Email request and/or writing-Email request to the transit server with the SMS mode. After performing a series of checks, such as an authentication check and an integrity check, the transit server forwards the requests to the Email-box server, and provides interactive services, which include a reading operation and a writing operation, for the mobile communication terminal and the Email-box server with the SMS mode.\nHowever, current interactive operation mode between the mobile communication terminal and the Email box has the following disadvantages.\nDue to the single mode for operating Email box, mobile communication terminal users may only complete Email-related operations by accessing an Email box with the SMS mode. The direct result brought about by the single operation mode is inconvenient operations, and low efficiency of interactive operations. Besides, when accessing the Email box with the SMS mode, the operation mode will be loaded down with trivial details, e.g., it is necessary to input instructions frequently, and there may be some limitations for length of the message content, thus efficiency of the interactive operations will be low."}
-{"text": "Photosensitive composition plays an important role in the fabrication of printed circuit boards (PCBs). Photosensitive compositions can be classified, based on their intended applications, into two main categories: photoresists for etching and solder resists. The former are used in the laying out of the copper circuit pattern using photolithographic techniques and will be removed after the copper circuit pattern is completed. The latter are used after the copper circuit pattern is laid out to protect certain portions of the copper circuit from being contacted by the soldering material during the soldering process. Both types of photosensitive compositions can be applied using either a liquid-type or a film-type technique. With the liquid-type technique, the photosensitive composition is applied onto the substrate using screen printing, roller coating, curtain coating, dip coating, etc. With the film-type technique, on the other hand, the photosensitive composition is first coated onto a transparent film, typically a transparent polyester film, which is then overlaid on the substrate and the photosensitive composition is then applied onto the substrate surface by hot-laminating or other suitable techniques.\nPhotosensitive materials, or photoresists in general, can also be categorized according to how they are developed. At the present time, there are two most common types of developments for photoresists: the solvent-development type and the alkaline-development type. The former uses a solvent, typically an organic solvent, to remove (for a negative type photoresist) the exposed portion of the photoresist. On comparison, if the photosensitive composition contains carboxyl groups, the exposed portions can be removed, and thus developed, using an alkaline solution such as sodium carbonate. Because the solvent-type development process can cause solvent exhaust and waste disposal problems which can cause significant increase in the overall manufacturing cost, the current industrial trend is to use the alkaline-development type processes. During the alkaline-development process, the carboxyl-containing polymer molecules will be neutralized by sodium carbonate to form a polymer salt. The unexposed portion of the photosensitive composition, which was not cross-linked, will be emulsified and caused to be dispersed in water. This thus completes the development process utilizing a negative-type photoresist.\nThe photosensitivity of a photosensitive composition is related to the degree of crosslinking after being exposed to UV light. Typically, a greater number of unsaturated double bonds (typically provided by the so-called photomonomers) in the photosensitive composition will result in a higher degree of cross-linking. However, excessive amounts of photomonomers would result in the photosensitive layer being too tacky, thus adversely affecting the image transfer process. With a dry film type photosensitive material, the photosensitive composition is protected by the transparent polyester film from being contaminated and the oxygen-induced polymerization-inhibition, and the requirement for the unsaturated double bonds can be entirely supplied by the photomonomers. Under these conditions, a alkaline-developable photoresist can be prepared by introducing only the carboxyl groups into the polymer binder, there is no need to add double bonds into the polymer binder.\nOn the other hand, when a solution-type photosensitive material is used, because no such protective polyester film is provided, the film thickness and the reactivity of the photosensitive film becomes important, and, as a result, the photosensitive composition can only contain a limited amount of photomonomers. In order to compensate for the decreased photosensitivity due to the reduction in the amount of photomonomers, unsaturated double bonds must be incorporated into the polymer binder so as to increase its post- UV exposure degree of cross-linking. This method has been disclosed in, for example, U.S. Pat. No. 5,055,378, which disclosed a photosensitive composition containing a polymer binder formed from the reaction of an epoxy resin with an carboxyl-containing unsaturated compound so as to impart unsaturated double bonds to the polymer binder. However, the polymer binder disclosed in the '378 patent does not contain carboxyl groups. As a result, it must be developed using a solvent and cannot be developed using an alkaline solution.\nU.S. Pat. Nos. 5,100,767, 5,009,982 and 4,943,516, and Japan Patents 61-243869 and 1-54390 disclosed various photosensitive compositions whose polymer binders are imparted with carboxyl groups by reacting the polymer binder disclosed in the '378 patent with anhydrides. Such an additional step not only allowed the photosensitive composition to become alkaline-developable, it also improved the photosensitivity of the polymer binder by introducing unsaturated double bonds into the polymer binder. However, because the formation of such alkaline-developable photosensitive polymer binder requires excess (about 10 to 20%) amounts of unsaturated carboxyl group containing compounds (such as acrylic acid or methylacrylic acid), which are typically volatile and can cause health concerns, the final finished products can cause health hazards to human bodies, especially after long term exposure to these products. This is one of the main disadvantages of the photosensitive compositions disclosed in the above-mentioned patents.\nU.S. Pat. Nos. 4,764,452 and 5,364,736 and Japanese patents 4-239070 and 2-47657 disclose photosensitive compositions containing polymer binders based on modified styrene-maleic anhydride resins, or the SMA resins. In these photosensitive compositions, partially-esterized SMA resins were reacted with saturated low-molecular weight alcohols in a ring-opening reaction by which carboxyl groups are introduced into the side chains of the SMA resins to form a carboxyl group containing polymer binder which consequently becomes base-developable. Because the reaction to prepare the base-developable polymer binder does not involve unsaturated carboxyl group containing molecules, it does not raise the health hazard problem mentioned above. However, these photosensitive compositions exhibited another type of shortcoming in that, because the polymer binder does not contain unsaturated double bonds, the developed, or hardened, final product does not provide adequate chemical resistance and hardness due to inadequate degree of crosslinking. In order to increase the degree of crosslinking for the post-exposure polymer binder, it has been proposed to add a mono-functional unsaturated alcohol, such as hydroxyethyl (meth)acrylate, in the ring-open reaction of the SMA resin. This method was disclosed in U.S. Pat. Nos. 5,087,552, 5,114,830 and 5,296,334. The use of mono-functional unsaturated alcohols not only imparts the carboxyl group into the SMA resin, it also imparts double bonds so as to increase the degree of cross-linking. However, while this action increases the photosensitivity of the polymer binder, it also returns the old problem of leaving low-molecular weight volatile, smelly and potentially health-hazard molecules in the polymer binder."}
-{"text": "As semiconductor designs continue marching towards smaller geometry, back end of line (BEOL) copper wire formation with lower dielectric constant material becomes a significant challenge. Damascene integration techniques are commonly used to form copper wires and vias. Conventional damascene integration methods involve forming trenches in an insulator layer, depositing a thin diffusion barrier layer and seed layer on surfaces of the trenches, electroplating to fill the trenches in the insulator layer, and planarizing, e.g., with chemical mechanical polish (CMP). As geometry becomes increasingly smaller, e.g., sub-micron, it is very difficult to sputter and plate copper material into tiny damascene trenches. Moreover, a significant amount of copper (Cu) wire thickness variation exists when Cu-CMP is employed."}
-{"text": "Conventionally, monitors enabling an observer to view two-parallax images captured from two viewpoints stereoscopically by using a specific device, such as a pair of stereoscopic vision glasses, have been in practical use. Furthermore, in recent years, monitors enabling an observer to view multi-parallax images (e.g., nine-parallax images) captured from a plurality of viewpoints stereoscopically with the naked eyes by using a beam control element, such as a lenticular lens, have also been in practical use. Such two-parallax images and nine-parallax images displayed on monitors enabling stereoscopic vision may be generated by estimating depth information of an image captured from one viewpoint and performing image processing with the information thus estimated.\nAs for medical image diagnosis apparatuses, such as X-ray computed tomography (CT) apparatuses, magnetic resonance imaging (MRI) apparatuses, and ultrasound diagnosis apparatuses, apparatuses capable of generating three-dimensional medical image data (hereinafter, referred to as volume data) have been in practical use. Conventionally, volume data generated by such a medical image diagnosis apparatus is converted into a two-dimensional image by various types of image processing, and is displayed two-dimensionally on a general-purpose monitor. For example, volume data generated by a medical image diagnosis apparatus is converted into a two-dimensional image that reflects three-dimensional information by volume rendering processing, and is displayed two-dimensionally on a general-purpose monitor.\nIn the conventional technology, however, when sequential multi-parallax images are generated and displayed in real time, processing load may increase, thereby making it difficult to display the sequential multi-parallax images smoothly."}
-{"text": "This invention relates generally to the surveying of static structures, such as buildings, and more particularly surveying such structures for real estate tax assessment purposes."}
-{"text": "The present invention concerns a method for conducting a process in an automatically controlled system. A mathematical model of the process, with at least one variable model parameter, is implemented in a computing means. Before the process is started, the model precomputes at least one selected process parameter according to input values supplied to it for presetting the system. The input values and the process parameter are measured during the process and the precomputed value of the process parameter is adaptively improved after the process based on the measured input values. Furthermore, the present invention concerns a device for implementing a method for conducting a process.\nA method and a device for conducting a quasi-continuous process in an automatically controlled system are discussed in German Patent Application No. 40 40 360. Such processes typically include rolling trains where each pass of a rolled strip forms a process cycle (hereinafter, \"process run\"). Like all actual industrial processes, these process runs are time-variable. To conduct such processes, the system controlling the process must be preset before each run. That is, unlike traditional closed-loop control, control must precede the actual process because in industrial processes, controlled values can often be measured only indirectly and not at the point of affection of the process. Consequently, direct closed-loop control is not possible.\nThe system controlling the process is preset in a known manner by precomputing selected process parameters according to pre-established input values, or initially estimated input values, or both, based on a pool of relevant mathematical models of the process and by presetting the system using the precomputed parameters. Since mathematical models of the process to be conducted can only approximate the actual process, the model must be adaptively matched to the actual process events. For this, the process parameters and the input values are measured directly or indirectly by determining other measured values during each process run. After the process run is complete, the precomputation performed with the mathematical models is repeated within the framework of a postcomputation done based on the input values subsequently measured. The variable model parameters are adaptively modified based on the deviation between the process parameters thus computed and the measured process parameters to reduce the deviations obtained. The model parameters thus adapted are available at the beginning of the following process run for precomputing the process parameters.\nWith very complex relationships, partial aspects of the process to be modeled are usually described in partial models, and the interaction among the partial aspects of the process is described by a higher order model that links the partial models. While the partial models are still capable of describing the respective partial aspects of the process with sufficient accuracy, model assumptions for linking the partial models are difficult to make and such models may be subject to serious errors. In particular, adapting the higher-order model in addition to adapting the partial models is very difficult because the partial models provide no exact measured values, but only estimated values as input values for the higher-order model.\nThe goal of the present invention is to improve the precomputed value of the process parameter in the case of a very complex relationship between the process parameter and the input values."}
-{"text": "1. Field of the Invention\nThe present invention relates to a method of measuring pyrophosphate.\n2. Description of the Related Art\nAs conventional techniques of measuring pyrophosphate enzymatically, there are known, for example, (i) a method involving using pyruvate orthophosphate dikinase and luciferase (Japanese Patent Application Laid-Open No. 2007-097471), (ii) a method involving using hypoxanthine phosphoribosyltransferase and xanthine dehydrogenase/oxidase (Japanese Patent Application Laid-Open No. 2003-174900), and (iii) a method involving using kinase, an enzyme for producing ATP from pyrophosphate, and dehydrogenase which requires NAD or NADP as a coenzyme (Japanese Patent Application Laid-Open No. 2006-187251)."}
-{"text": "This invention relates to air metering apparatus and more particularly to such apparatus for use with means for controlling the ratio of air to fuel in a mixture to be combusted in an internal combustion engine.\nThe control of emissions from internal combustion engines and particularly automobile engines has become a major environmental concern. Various federal and state regulatory agencies have promulgated emission standards for certain substances found in the combustion products entering the atmosphere through an engine's exhaust, the most important of these substances being hydrocarbons, carbon monoxide and oxides of nitrogen. To meet emission control standards, various pollution control devices such as catalytic converters and thermal reactors have been developed for use with automobile engines to reduce the quantities of unwanted substances emitted into the atmosphere to within prescribed limits.\nIt has been found that most efficient removal of unwanted substances by pollution control devices is achieved when an engine is operated within a narrow range of air-fuel ratio values for an air-fuel mixture combusted in an engine. Consequently, numerous systems have been developed which attempt to maintain the air-fuel ratio of a mixture to be combusted in an engine within this value range. Examples of systems of this type are disclosed in U.S. Pat. Nos. 3,939,654, 3,946,198, 3,949,551 and 3,963,009. While the systems disclosed in these patents do tend to keep the air-fuel ratio for a mixture to be combusted within the value range where maximum efficiency in removal is obtained, this is usually accomplished only by constantly adjusting the air-fuel ratio. Further, overadjustments frequently occur which then require additional corrections and the systems respond to transitory changes in an engine's operating characteristic to make adjustments when none are actually needed."}
-{"text": "Hovercrafts are vehicles that travel across land or water just above a cushion of air. Generally, a skirt extends down from the periphery of the vehicle to the land or water in order to entrap the supporting cushion of air. Depending on the size of the vehicle and/or its intended application, it may be necessary to provide a stability system that reduces the effects of roll and pitch forces. For example, a longitudinally extending keel system can be attached to the underside of the vehicle hull along the vehicle's center line. Ideally, the keel stays in contact with the water's surface to prevent air on one side thereof from moving to the other side thereof. In this way, differential pressure cells are created to provide roll stability. Similarly, a laterally extending keel system can be used to provide pitch stability.\nOne such keel system is utilized by the U.S. Navy in its LCAC (landing craft, air cushion) vehicle. As shown in the plan view of FIG. 1, the LCAC has a segmented longitudinal keel system 20 and a segmented lateral keel system 50 contained within the periphery of a hovercraft's skirt (not shown). Longitudinal keel system 20 extends along the direction of craft motion which is indicated by arrow 100. Lateral keel system 50 is transverse to longitudinal keel system 20 and consists of a plurality of inflated cones 52.\nOne of cones 52 of lateral keel system 50 is shown in a side view in FIG. 2 which is a view taken along line 2--2 of FIG. 1. The side view in FIG. 2 depicts a triangle (of height h and base dimension b) that is essentially positioned parallel to craft motion 100. The portion of cone 52 parallel to the load forces determines the cone's ability to resist the buckling forces associated with craft motion 100. This is best understood by examining the relationship used to determine bending stress .sigma..sub.bending, i.e., the force that causes buckling of a cone 52 subjected to craft motion 100. Bending stress .sigma..sub.bending can be defined mathematically as ##EQU1## where M is the bending moment (i.e., the applied force at the cone's tip or apex times the height h of the cone), c is the distance from the outermost surface of a cone to its central axis (i.e., in this case, one half of the cone's width w as indicated in FIG. 1), and I is the moment of inertia. The moment of inertia I for a triangle is defined as ##EQU2## Since .sigma..sub.bending and I are inversely proportional, the b and h values of each cone 52 must be sufficient to keep .sigma..sub.bending from exceeding a threshold at which the (inflated) cone will buckle.\nThe existing longitudinal keel system is shown in greater detail in isolation in FIG. 3 and in cross-section in FIG. 4 which is taken along line 4--4 of FIG. 3. The longitudinal keel system is attached to the underside of the LCAC's hull which is shown in part and referenced by numeral 10. The longitudinal keel system consists of upper inflatable bladder sections 20A, 20B and 20C (extending fore to aft) attached to the underside of hull 10. Each upper inflatable bladder section 20A, 20B and 20C is trough-shaped and is inflated or pressurized by the craft's lift/fan system (not shown). Attached to each section 20A, 20B and 20C is a corresponding lower trough-shaped inflatable bladder section 22A, 22B and 22C. Inflation of sections 22A, 22B and 22C is accomplished by the provision of a plurality of inflation holes 24 (see FIG. 4) passing between the upper and lower inflatable lower bladder sections 20A and 22A. Thus, inflation of bladder sections 22A, 22B and 22C is accomplished when air passes through these holes. Port and starboard longitudinally extending blades 26 and 28 are attached to the lower portion of sections 22A, 22B and 22C. Blades 26 and 28 are flexible and oppose one another to define a channel therebetween. Holes 30 between blades 26 and 28 in the lower portion of lower inflatable bladder sections 22A, 22B and 22C provide an air exhaust and allow for the removal of water that might seep into the inflated sections while blades 26 and 28 cooperate to prevent water from entering holes 30.\nAlthough accomplishing its intended function of roll stability, the longitudinal keel system shown and described above has been susceptible to environmental abrasion and impact damage that propagates forward and aft from the point of damage. Since the aft-most lower inflatable bladder section 22C is the section most likely to be damaged, such damage tends to propagate forward to section 22B. Damage greatly affects the craft's roll stability and is further costly to repair owing to the location of the keel system."}
-{"text": "Solar power or solar energy has been used for many decades for heating dwellings and water and for generating electricity. Because solar power is a renewable energy source much effort has been made to develop systems to use such energy. The costs have been high and the storage of energy has not been very effective. A significant problem is the need for effectively positioning solar concentrators relative to the sun as the sun moves during the day. The best thick plastic-fibers available in the market have been studied but none of these provided the required level of optical transmission in the near-infrared (NIR) and a large fraction of the incident energy is lost in the waveguides after only a few meters propagation. The efficiencies of conventional systems have been low and there is a need for a more efficient and cost-effective system.\nThe method of the present invention provides a solution to the above-outlined problems. More particularly, the present invention relates to a solar power system that has a solar concentrator that has a lens disposed therein. The solar concentrator has a first cone-shape and a bottom. The solar concentrator has a length (l) that is longer than a width (w) at the lens of the solar concentrator. A solid tapering device is disposed at the bottom and has a second cone-shape. The tapering device has a bottom in operative engagement with a first curved glass loop section of a glass rod. The glass rod has a second curved glass loop section and a straight section. The first curved loop section and the second curved glass loop section have a first gap defined therebetween. The second curved glass loop section and the straight glass section have a second gap defined therebetween. The solar concentrator is adapted to receive solar power as rays and conveys and concentrates rays as light towards a focus line or segment at the tapering device. The tapering device is in communication with an upper end of the first curved glass loop section. The first curved glass loop section conveys the light to the second curved glass loop section via and across the first gap. The second curved glass loop section conveys the light to the straight glass section via and across the second gap. While transmitting light, the first curved glass loop section is rotated relative to the second curved glass loop section at the first gap. The second curved glass loop section is rotated relative to the straight glass section at the second gap so that the solar concentrator follows a path of a sun. An outer end of the straight glass section is disposed in proximity to a water surface to heat the water below the water surface.\nIn an alternative embodiment, the outer end is disposed above the water surface.\nIn yet an alternative embodiment, the outer end is disposed above a water surface of water contained in a container and the water surface is maintained between a minimum level and a maximum level.\nIn an alternative embodiment, the light emitted from the outer end converts the water to steam that is conveyed from the container."}
-{"text": "This invention relates to produce labelers and, in particular, to selectively labeling a wide variety of fresh produce by size or grade at high speed.\nWithin the fruit packing industry it has become common practice for fruit to be labeled in large quantities at high speed on sizing/grading equipment. The demand for labeling is driven mainly by the retail industry, requiring product to be identified by variety and/or source and furthermore by size and quality. As a consequence, packers are commonly required to apply a number of different labels to fresh produce during the grading and packing operation. This is generally achieved by installing multiple labeling stations in sequence over the grading conveyors. Two major drawbacks exist with this arrangement: (a) the extended space required to accommodate multiple labeling stations, and (b) the resulting higher costs attached to such installations.\nFor example, if peaches are to be labeled \u201csmall, medium or large,\u201d the peaches typically pass through sizing equipment where three banks of labeling equipment are used to apply the appropriate size labels to the sized peaches. The 3 labeling stations (including rotary bellows, etc.) all take up space and are all relatively expensive.\nThere is a clear need for labeling equipment that is more efficient and versatile than the prior art systems that use separate labeling stations for each separate size or grade of produce.\nThe closest prior art known to applicant includes the use of dual cassettes as shown in Rietheimer U.S. Pat. No. 5,645,680 (see FIG. 6). However, Rietheimer has three major weaknesses. First, he uses a complex fixed cam surface housing (as opposed to rotary bellows). Secondly, his system uses conventional knife edge label stripping and is therefore limited to the use of relatively stiff labels. Thirdly, Rietheimer requires that the guide plate of each label cassette discharge each label tangentially to the rotary applicator and parallel to each transfer head as shown in FIG. 6 herein. This geometry limits the number of cassettes usable to two (see FIG. 6). The present invention is capable of using three or more label cassettes without the complex camming mechanism of Rietheimer. The present invention also uses thin and flexible labels, which are usable on more items than stiff labels.\nThe present invention provides, for the first time, a single automatic labeling station capable of applying 3 or more different labels to singulated produce passing through the station. The present invention also provides, for the first time, a single automatic labeling station wherein 3 or more label cassettes interact with a single rotary bellows applicator.\nIn the above example of \u201csmall, medium and large\u201d peaches to be labeled, the present invention labels all 3 separate sizes in a single station with a single rotary bellows applicator. The invention reduces most of the prior art machinery required and the space necessary to house the machinery!\nA primary object is to provide an automatic produce labeling apparatus capable of applying 3 or more different labels at a single labeling station and at high speed, i.e., more than 1000 labels per minute.\nA further object is to provide an automatic labeling system which eliminates the need for multiple labeling stations otherwise required by the prior art in applying labels displaying different sizes or grades on produce.\nA further object is to provide labeling apparatus wherein multiple label cassettes interact with a single rotary bellows applicator to apply multiple different labels to produce at speeds in excess of 1000 labels per minute.\nFurther objects and advantages will become apparent from the following description and drawings wherein:"}
-{"text": "Devices such as laser diodes are used in various different applications. Many applications require the laser diode to be on for a relatively short period of time. If a device such as a laser diode is on for a prolonged length of time, due to a fault, there may be a risk that safety is compromised. For example, such devices are often used in consumer products such as mobile phones. Safety may be a particular concern in such devices. Alternatively or additionally the problem may be that if a device is on when it should be off, excess power is consumed. This may be disadvantageous in situations where power is provided, for example, by a battery."}
-{"text": "A broadcast network uses the Moving Picture Experts Group-2 Transport Stream (MPEG-2 TS) for delivering multimedia content according to the related art. The MPEG-2 TS is used as a representative transmission technique for delivering a bit stream in which a plurality of broadcast programs (a plurality of encoded video bit streams) can be transmitted without having an error. For example, the MPEG-2 TS is appropriate for digital TV broadcasting, or the like, in the multimedia era.\nFIG. 1 illustrates a layer structure for supporting the MPEG-2 TS according to the related art.\nReferring to FIG. 1, the layer for supporting the MPEG-2 TS includes a media coding layer 110, a sync layer 120, a delivery layer 130, a network layer 140, a data link layer 150, and a physical layer 160.\nThe media coding layer 110 and the sync layer 120 are formed of a format which is easily usable as a basic unit of recording or transmission of media data. Further, the delivery layer 130, the network layer 140, the data link layer 150, and the physical layer 160 record a data block (for example, \u201cAU\u201d) configured by the sync layer 120 in a recording medium or configure a multimedia frame for transmission. The configured multimedia frame is delivered to a subscriber's terminal, or the like, through a predetermined network.\nThe sync layer 120 is composed of a fragment block 122 and an access unit 124, and the delivery layer 130 is composed of an MPEG-2 TS/MP4 132, an RTP/HTTP 134, and a UDP/TCP 136.\nHowever, the MPEG-2 TS has several limits in supporting a multimedia service. That is, the MPEG-2 TS has the limitations of unidirectional communication, non-efficiency in transmission due to a fixed frame size, and a generation of an unnecessary overhead while delivering data by using a transmission protocol and an Internet Protocol (IP) specialized to an audio/video.\nThe MPEG newly suggests the MPEG Media Transport (MMT) standard as one of the multimedia transmission techniques for supporting the multimedia service based on the MPEG technique. The MMT standard is suggested by the MPEG in order to overcome the limits of the MPEG-2 TS.\nThe MMT standard may be applied in order to efficiently deliver hybrid contents through heterogeneous networks. Here, the hybrid contents refer to a set of contents having combined multimedia elements, such as a video/audio/application, or the like. The heterogeneous network means a network in which a broadcast network, a communication network, or the like, are mixed.\nIn addition, the MMT standard defines a transmission technique more friendly to the Internet Protocol (IP) which serves as a basic technique for transmitting the multimedia.\nThe MMT standard typically proves to be an efficient MPEG transmission technique in a multimedia service environment, varying based on the IP; continuous research and standardization thereof have progressed.\nThe MMT standard requires a preparation of a method of providing an efficient MPEG transmission technique in a recent multimedia service environment aiming to provide a hybrid network and hybrid contents.\nTherefore, a need exists for a method and apparatus for configuring a control message to be used in a broadcast system that supports an IP-based multimedia service. The above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
-{"text": "1. Field of the Invention\nThis invention relates to InGaN/GaN light emitting diodes (LEDs) and laser diodes (LDs), and more particularly to nonpolar III-nitride LEDs in which the wavelength can be controlled by selecting the barrier thickness of multiple quantum wells (MQWs).\n2. Description of the Related Art\nThe usefulness of gallium nitride (GaN) and its ternary and quaternary compounds incorporating aluminum and indium (AlGaN, InGaN, AlInGaN) has been well established for fabrication of visible and ultraviolet optoelectronic devices and high-power electronic devices. These compounds are referred to herein as Group III nitrides, or III-nitrides, or just nitrides, or by the nomenclature (Al, B, Ga, In)N. Devices made from these compounds are typically grown epitaxially using growth techniques including molecular beam epitaxy (MBE), metalorganic chemical vapor deposition (MOCVD), and hydride vapor phase epitaxy (HVPE).\nCurrent nitride technology for electronic and optoelectronic devices employs nitride films grown along the polar c-direction. However, conventional c-plane quantum well structures in III-nitride based optoelectronic and electronic devices suffer from the undesirable quantum-confined Stark effect (QCSE), due to the existence of strong piezoelectric and spontaneous polarizations. The strong built-in electric fields along the c-direction cause spatial separation of electrons and holes that in turn give rise to restricted carrier recombination efficiency, reduced oscillator strength, and red-shifted emission.\nOne approach to eliminating the spontaneous and piezoelectric polarization effects in Group-III nitride optoelectronic devices is to grow the devices on nonpolar planes of the crystal. For example, in GaN crystals, such planes contain equal numbers of Ga and N atoms and are charge-neutral. Furthermore, subsequent nonpolar layers are equivalent to one another so the bulk crystal will not be polarized along the growth direction. Two such families of symmetry-equivalent nonpolar planes in GaN are the {11-20} family, known collectively as a-planes, and the {10-10} family, known collectively as m-planes.\nThe other cause of polarization is piezoelectric polarization. This occurs when the material experiences a compressive or tensile strain, as can occur when (Al, B, Ga, In)N layers of dissimilar composition (and therefore different lattice constants) are grown in a nitride heterostructure. For example, a thin AlGaN layer on a GaN template will have in-plane tensile strain, and a thin InGaN layer on a GaN template will have in-plane compressive strain, both due to lattice matching to the GaN. Therefore, for an InGaN quantum well on GaN, the piezoelectric polarization will point in the opposite direction than that of the spontaneous polarization of the InGaN and GaN. For an AlGaN layer lattice matched to GaN, the piezoelectric polarization will point in the same direction as that of the spontaneous polarization of the AlGaN and GaN.\nThe advantage of using nonpolar planes over c-plane nitrides is that the total polarization will be reduced. There may even be zero polarization for specific alloy compositions on specific planes. Such scenarios will be discussed in detail in future scientific papers. The important point is that the polarization will be reduced compared to that of c-plane nitride structures.\nAlthough high performance optoelectronic devices on nonpolar on-axis m-plane GaN have been demonstrated, it is difficult to obtain long wavelength emission from InGaN/GaN MQWs grown on m-plane GaN. This is probably due to the low In incorporation of the InGaN/GaN MQWs. The emission wavelength of devices grown on m-plane is typically 400 nm, while the emission wavelength of devices grown on c-plane is 450 nm at the same growth condition. Reducing the growth temperature increases the In incorporation; however, crystal quality would be degraded. This would be a significant problem for applications such as blue, green, yellow, and white LEDs.\nThe present invention describes a technique for the growth of nonpolar III-nitride light emitting devices in which the emission wavelength from the devices can be controlled by the barrier thickness of the MQWs in the devices. For example, the present invention has obtained blue and green emission without the effect of polarization."}
-{"text": "Air distillation plants generally include adsorbers so as to subject the air, compressed on entering the plant, to a drying and decarbonation treatment before it is cooled and then introduced into the distillation column proper. For this purpose, the air to be treated flows under a pressure typically of about 5 bar absolute through a bed of an adsorbent based in the adsorber. Once the adsorption capacity of the bed has been exhausted, the adsorbent is regenerated by means of a dry and decarbonated gas. This gas may be an off-gas recovered from the distillation column and injected into the adsorber in the reverse direction. In order to increase the regeneration effectiveness, this off-gas is heated to a high temperature of about 200.degree. C. This type of adsorber, in which the adsorbent is subject to cyclic temperature variations, is known by the name \"TSA (Temperature Swing Adsorption) adsorber\".\nHeating the gas as well as the adsorbent represents a major part of the energy consumed by the distillation plant. In order to minimize this consumption, it is therefore necessary to reduce the heat exchange between the adsorber and its environment.\nThermal insulation means consisting of an outer covering of the walls of the adsorber are known. This type of insulation proves to be very effective for the entire adsorber but does have certain drawbacks when the adsorber is subjected in a cyclic manner to temperature variations. This is because, since the walls of the adsorbers are generally made of metal, these constitute, depending on the phase of the cycle a heat store or heat sink tending to delay establishment of a steady state after a temperature change. In the regeneration phase, the walls absorb part of the heat provided by the heated gas and thus delay the heating of the adsorbent, in particular in regions close to the walls. In the adsorption phase, they restore this heat, thereby decreasing the adsorption effectiveness of the adsorbent. The walls therefore act out of step with the cycles and considerably decrease the effectiveness of the adsorber.\nThis phenomenon could be counteracted by increasing the size of the adsorbers and consequently by reducing, for a constant flow, the frequency of the cycles. However, recent developments in distillation plants have led to a reduction in the size of the adsorbers so as to reduce the amount of adsorbent and the cost of an adsorber. Because of the fact that the amount of adsorbent is reduced, the frequency of the cycles is considerably increased.\nIn order to improve the effectiveness of the known absorbers, adsorbers are known which are equipped with internal thermal insulation. Such internal thermal insulation must satisfy particular requirements:\nit must not allow bypassing of gas between the wall of the adsorber and the bed of adsorbent, so that the gas that is to be treated and the gas that has been treated are always separate; PA1 it must withstand the repeated compressions and decompressions imposed by the operating cycle and be able to attain pressure variations of about ten bar; and PA1 it must easily allow periodic hydrostatic testing, this being generally required by the authorities in the operating countries. PA1 the blocking means include at least one separating wall fitted between the said wall and the shell; PA1 the adsorber comprises means for decreasing natural convection in the gas-filled cavity; PA1 the means for decreasing natural convection in the gas-filled cavity comprise a sheet, in particular a metal sheet, placed between the wall and the shell in the gas-filled cavity in order to obtain a double gas-filled cavity; PA1 the sheet is centered with respect to the said wall and to the shell; PA1 the gaps between the sheet and the wall and between the sheet and the shell are between 10 and 40 mm, and preferably equal to 20 mm; PA1 the absorber comprises elements, especially point or linear elements, forming spacers between, on the one hand, the wall and the sheet and, on the other hand, between the sheet and the shell; PA1 the elements forming spacers between the wall and the sheet are offset with respect to the elements forming spacers between the sheet and the shell; PA1 the elements forming spacers are indented portions and raised portions preferably produced by goffering in the sheet, these indented portions and raised portions having, in particular, pointed ends; PA1 the indented portions and raised portions of the sheet extend parallel to the axis of the shell, over the entire length of that dimension of the sheet which is parallel to this axis, and have a cross-section of approximately triangular shape; PA1 the indented portions and raised portions alternate on each face of the sheet; PA1 in the case of an adsorber having a shell which is extended into a part forming a dome, the adsorber includes at least one and preferably two layers of a felt of fibres applied against the inside of the dome; PA1 for an adsorber intended, in particular, to operate under a varying pressure, the felt of fibres includes ceramic fibres and has a ratio between the fluid resistivity and Young's modulus of elasticity of less than 1 s/m.sup.2 ; PA1 the thickness of the or each layer of felt is between 20 and 30 mm, and preferably equal to about 25 mm; PA1 the layers of felt are attached to the inside of the dome by means of removable attachment devices; PA1 the removable attachment devices comprise rods, fastened to the dome and intended to transpierce the layers of felt, and washers for holding the felt against the dome, the washers engaging with the rods, in particular by a clip-in mechanism; and PA1 the adsorber is used in particular for the purification of air sent into an air distillation plant, the active material being an adsorbent.\nDocument US-A-3,925,041 teaches a TSA adsorber equipped with internal thermal insulation. This adsorber comprises a cylindrical shell oriented vertically, the ends of which are closed by ellipsoidal parts forming, respectively, a dome and a bottom. A bed of adsorbents is placed in the cylindrical shell\nIn the region of the bed of adsorbent, the internal thermal insulation is produced by sheets of a rigid insulant, these being applied end to end against the shell and compressing, against this shell, a layer of a fibrous insulating material. This thermal insulation is therefore placed between the shell and the bed of adsorbent. In order to prevent bypassing of the gases through this insulation during operation, metal sheets are applied against the shell, covering the sides of the sheets of insulant so as to deflect the flow of gas towards the bed of adsorbent.\nFurthermore, a fibrous insulant is applied by means of bolts against the internal parts of the dome and of the bottom of the adsorber.\nAlthough this internal thermal insulation considerably increases the effectiveness of an adsorber, it is not suitable for undergoing a hydrostatic test. This is because, after such a test, the thermal insulation is wet and ineffective. It must therefore undergo lengthy, difficult and consequently very expensive drying. Moreover, fitting the sheets of insulant and the metal sheets in the region of the bed, of adsorbent proves to be difficult and lengthy so that only specialists are able to place this insulation against the internal walls of the adsorber. This known insulation therefore considerably increases the cost of manufacturing of the adsorber."}
-{"text": "1. Field of the Invention\nThe pesent invention relates to a supervisory device for wrong operation of varying quantity measuring apparatus, and more particularly to a supervisory device for wrong operation of varying quantity measuring apparatus which measures a physical varying quantity by relative movement between a main scale and an index scale.\n2. Description of the Prior Art\nThere has been a well-known physical varying quantity measuring apparatus which measures a certain physical varying quantity by relative movement between a main scale and an index scale, and is of wide use for the measurement of length and physical varying quantity in length measuring instruments, coordinate measuring machines, etc. or in positioning or the like of machine tools, or for the other purposes.\nThis varying quantity measuring apparatus includes a main scale and an index scale which are moved relative to each other in accordance with a physical varying quantity and an encoder which detects the relative movement of such two scales and measures the requested physical varying quantity based upon the relative movement of such two scales which can be detected by this encoder.\nThe encoder mentioned above detects the relative movement of both two scales as an electric signal by optical, electromagnetic or well-known means using contact resistance, and further converts the detected signal into a count pulse signal in accordance with the relative movement per unit time to supply this to a counter provided in the encoder. The counter counts the supplied count pulse signal, and the counted value of the count pulse signal is displayed as a digital value representing a measured value of the physical varying quantity or is used for an other control signal.\nAccordingly, this varying quantity measuring apparatus can measure the physical varying quantity with accuracy as a predetermined digital signal by means of relatively moving the main scale and the index scale in accordance with the desired physical varying quantity.\nThe varying quantity measuring apparatus mentioned hereinabove, however, clears the counted value of the count pulse signal stored in the encoder, when the electric power stops or lowers the voltage to great extent for some reason, and restart counting the count pulse signal from the very beginning, when the electric voltage recovers up to the optimum level. Accordingly, when the electric power stops or lowers its voltage, there arise measuring errors between the relative moving quantities of the main scale and the index scale which are converted into the electric signal by the encoder to be displayed and the actual moving quantities of these two scales. Furthermore, it is a problem that the error occurrence is not known to those at measuring work in most cases.\nEspecially, in such a case that the electric power suspension or low voltage is instantly caused by the start or the stop of a large machine, the sudden load change, or the like for example, the encoder immediately restart counting the count pulse right after the counted value of the count pulse signal is cleared, which is not known to those at measuring work in most cases, and an effective countermeasure has been requested."}
-{"text": "1. Field of the Invention\nThe present invention relates to an intake and exhaust gas claning apparatus, and more particularly to an exhaust gas cleaner for automobiles, motorcycles and factories.\n2. Description of the Prior Art\nAir pollution is becoming a serious problem worldwide. Engine exhaust of automobiles and motorcycles, and factory discharges are the major sources of pollutants. Engine exhaust pollutants of a motor vehicle are generated mostly by the following causes:\n1. The low temperature in the combustion chamber when the engine is starting from cold. PA1 2. Driving uphill, bumpy road conditions or heavy loads. PA1 3. Poor quality of gasoline or oils. PA1 4. Clogged air filter. PA1 5. Intake of poor quality air from a tunnel or congested road. PA1 6. Engine abrasion or too many pollutants in the combustion chamber. PA1 7. Back pressure caused by a smoke filter or catalytic converter. PA1 1. High volume and high temperature. PA1 2. High speed and impulsive vibration. PA1 3. Viscosity after dispersion. PA1 4. High vibration frequency and noise level. PA1 5. Contains fine particles and becomes greasy. PA1 6. Contains water vapor. PA1 7. Contains gases of carbon dioxide, sulfer dioxide, carbon monoxide, etc. PA1 8. Contains polynuclear aromatic hydrocarbons (PAH) and other cancer inducing or cell mutation gases. PA1 9. Indirectly forms peroxyacety/nitrate (PAN) and other photochemicals which are environmentally harmful. PA1 1. Wet type dust collector: tends to be blown away by vapor, consumes much water, tends to form erosive gas, not economical or pratical. PA1 2. Honeycomb filter gauze formed of ceramic or metal: easily clogged by greasy smoke, incur back pressure, difficult operation, easily damaged, high cost, time consuming, short effective life, requires considerable backup materials. PA1 3. Whirlwind type or static collector: large size, consumes much electric power, heavy load, not suitable for motor vehicles. PA1 4. Additives in gasoline or diesel oil: while there are claims that the additives can improve combustion, not much evidence or proof is available for their effectiveness in eliminating toxic gas, and not conclusive as to whether or not the additives will be harmful to the engine. PA1 5. Combustion smoke filter: provides better results, but high cost, consumes energy and oxygen, and produces more carbon dioxide gas, has explosion risk, limited effective life, not recyclable.\nEngine exhaust gas includes smoke which contains solid particles and liquid drops. They all create severe air pollution problems and threaten the health of human beings and animals. Some of the exhaust gases discharge into the atmosphere and form \"acid rain\" which further creates a water pollution problem. The carbon dioxide in engine exhaust gas is also one of the factors causing the global greenhouse effect. Thus, to minimize engine exhaust pollution is becoming a worldwide focus.\nEngine exhaust gases have the following characteristics:\nMany types of engine exhaust controlling or eliminating means have been developed, but they have disadvantages, as indicated below:\nThe above disadvantages have limited conventional exhaust controlling devices from being widely used. Thus, there has been a need for providing a more cost effective engine exhaust controlling means."}
-{"text": "The present invention relates to an improved drum assembly for washing, drying machines, and the like, as well as a method for making said drum assembly.\nThe drum assemblies of the type concerned by the present invention conventionally comprise a body of substantially drum configuration. This drum body is formed by a perforated metal band constituting the cylindric wall of the drum body, by a rear closure disc for closing the drum body, which supports the cross member supporting the shaft of the drum, and a front ring element defining the front loading mouth therethrough the cloth articles to be processed are loaded into the drum. In the case of a top loading, on the other hand, the front ring element is replaced by a closing cover, whereas the side band is provided with an opening for introducing the cloth articles into the drum.\nIn order to provide an efficient tumbling of the cloth articles to be washed or dried, on the inside surface of the perforated metal band forming the cylindric wall of the drum are provided a plurality of projections, having a substantially wedge-like cross section, which operate as xe2x80x9centrainmentxe2x80x9d elements. Actually, owing to the provision of these projections, the rotary movement applied to the drum assembly will cause the linen loaded therein to be upwardly entrained and fall by gravity, thereby aiding the washing operation.\nThe provision of the above mentioned projections, which will be hereinafter called xe2x80x9cwedge elementsxe2x80x9d, causes, however, problems or difficulties in properly coupling the metal band and ring element or front cover. Because of the provision of the mentioned wedge elements, in fact, the edge or rim of the mentioned metal band will have an uneven contour, which, on the other hand, must be perfectly fitted or mated with the contour of the ring element or front cover of the drum.\nPrior solutions for providing such a connection include the use of spot welding operations between the respective walls of the ring or front cover element and the drum wedge elements, formed on the corresponding surface engagement portions. Such a spot welding method, however, could alter or modify the required chemical-physical properties of the drum material (chromium stainless steel, AISI 430), which is would be seriously affected by a thermal treatment like that deriving from a spot welding operation. On the other hand, this spot welding operation, which must involve a number as low as possible of regions of the drum, does not provide a connection having the desired strength and duration. In fact, the welding spots do not provide a structurally continuous connection and, moreover, as stated, could undesirably modify the drum material thereby favoring oxidation phenomena which would be very dangerous for the time stability of the performed connection.\nIn addition to the above mentioned drawbacks, to weld the front ring element or cover on the wedge elements of the drum is disadvantageous since it does not allow a proper centering of these elements. Actually, as the welding operation is performed, the ring or cover element can be displaced from a proper position centered with respect to the cylindric band on which they must be locked, thereby the end product may be affected by openings or undesired gaps between the connected portions.\nTo the foregoing it should be added that for welding the drum, additional welding materials and power are required, thereby increasing the end cost of the drim assembly.\nAccordingly, the aim of the present invention is to provide such a drum assembly for washing, drying and the like machines which has a mechanical strength greater than that of commercially available like devices, mainly at the level of the mutual connection regions coupling the mentioned wedge elements and front ring or cover element.\nWithin the scope of the above mentioned aim, a main object of the present invention is to provide a drum assembly of the above mentioned type which is effectively devoid of any unevenness like those which occur in the connection regions of the wedge elements and the front ring or cover element in prior drum assemblies.\nYet another object of the present invention is to provide a method for making the above mentioned drum assembly, which is specifically designed for reducing to a minimum the end cost of washing, drying and the like machine drum assemblies.\nAccording to one aspect of the present invention, the above mentioned aim and objects are achieved by an improved drum assembly for washing, drying and the like machines, of the type including a cylindric band provided with a plurality of wedge elements facing the inside of the drum assembly, a closing rear disc and a front ring element, characterized in that said drum assembly comprises connecting means for directly mechanically connecting said ring element and band, at said wedge elements of said band.\nAccording to further features of the inventive drum assembly, said connecting means comprise projecting or raised walls provided on that section of each said wedge element engaging said ring element, said projection portions being adapted to cooperate with corresponding projecting or raised portions of said ring element. Said projecting portions, in particular, comprise a wall formed on the end of the wedge elements facing the ring element and directed inwardly of said wedge elements, at a position substantially parallel to the surface of the ring element therewith said wall must be connected. In turn, said projecting portions of said ring element comprise a deformed or upset edge or rim portion provided at the connection section on each wedge element and having such a size and shape to be housed inside each said wedge element, above and beyond the mentioned projecting wall.\nAccording to further features of the drum assembly according to the present invention, said deformed edge or rim portion comprises a wall facing the inside of the drum assembly. After having performed the mechanical connection, the mentioned deformed edge or rim portion will be rearwardly turned and upset against said wall of the respective wedge element.\nAccording to a modified embodiment of the drum assembly of the invention, said direct connection mechanical means comprise a portion of said wedge elements axially projecting on the outside of said cylindric band, for cooperating with a corresponding window of the front ring element, to be locked inside the latter, at the contour thereof.\nThe method for making the above disclosed drum assembly, which also constitutes a main aspect of the present invention, is essentially characterized in that it comprises the step of cold anchoring said front ring or cover element and cylindric band, at the wedge elements of said band. The method is moreover characterized in that said anchoring or connecting step is performed by an upsetting operation which is preceded by a self-centered connection between said front ring or cover element and said cylindric band of the drum assembly.\nFurther features of the inventive method are clearly defined in the remaining dependent claims.\nWith respect to prior drum assemblies, the inventive drum assembly provides the advantage of allowing a mechanical direct connection of the wedge elements and front ring or cover element, i.e. a connection which does not comprise any welding spots favouring oxidation and sealing failure phenomena of the performed connection.\nThe mechanical sealing properties of the subject connection are moreover much more stable and efficient, than prior welding spot connections, since the subject connection is extended to the overall mutual coupling section of the wedge elements and corresponding portion of the front ring or cover element, and not to only a portion thereof, as in a spot welded connection.\nThe cold connection performed by upsetting a projecting portion of said front ring or cover element on the contour of a corresponding wedge element, furthermore, provides a self-centering effect. Thus, the drum assembling operation can be carried out a fully automatic manner due to the mutual engagement of the portions being connected, thereby providing a well finished final product."}
-{"text": "The term powertrain, or powerplant, describes the main components that generate power and deliver power to the road surface, water, or air. In a wider sense, the powertrain includes components used to transform a stored (chemical, solar, nuclear, kinetic, potential, etc.) energy into kinetic energy for propulsion purposes.\nThe powertrain of the current and upcoming electric vehicles (EVs) and plug-in hybrid electric vehicles (PHEVs), cumulatively called plug-in electric vehicles (PEVs), is composed of a high-power battery pack, as the Energy Storage Sub-System (ESS), and an Inverter followed by a Propulsion Machine for electric propulsion.\nElectric vehicles are vehicles propelled by electricity as opposed to the conventional vehicles which operate by liquid or gaseous fuels. Electric vehicles may use a combination of different energy sources such as, for example, fuel cells (FCs), batteries, and super-capacitors (or ultra-capacitors) to power an electric drive system as shown in FIG. 1. Electric vehicles (EVs) operate with high energy conversion efficiency, produce a lower level of exhaust emissions, and lower level of acoustic noise and vibration than conventional vehicles. The electricity needed for electric vehicles' operation may be produced either external the vehicle and stored in the battery, or produced onboard with the help of fuel cells.\nA plug-in electric vehicle (PEV) is any motor vehicle that can be recharged from an external source of electricity, such as wall sockets. The electricity stored in the rechargeable battery packs drives or contributes to drive the wheels. PEV is a superset of electric vehicles that includes all-electric or battery electric vehicles (BEVs), plug-in hybrid vehicles (PHEVs), and electric vehicle conversions of hybrid electric vehicles as well as conventional internal combustion engine vehicles.\nIn electric vehicles, the main energy source is assisted by one or more energy storage devices. A combination of batteries and super-capacitors are often used as Energy Storage Sources which can be connected to the fuel cell stack in a number of ways. The voltage characteristics of the two devices must match perfectly, and only a fraction of the range of operation of the Energy Storage devices can be utilized. For example, in a fuel cell/battery configuration, the fuel cell must provide substantially constant power all the time due to the fixed voltage of the battery, and in a battery/super-capacitor configuration, only a fraction of the energy exchange capability of the super-capacitor (also referred to herein as ultra-capacitor) can be used.\nDC-to-DC converters can be used to interface the elements in the electric powertrain by boosting or chopping the voltage levels as required by the load in different regimes of the EV operations. By introducing the DC-to-DC converters, one can select the voltage variations of the devices, and the power of each device can be controlled.\nDifferent configurations of EV power supplies dictates that at least one DC-to-DC converter is necessary to interface the fuel cell, the battery, or the super-capacitor's module and the DC Link. The DC-to-DC converter structure has to be reliable and lightweight, small in volume, operate with high efficiency, and low electromagnetic interference, and low current/voltage ripple.\nA DC-to-DC converter is a category of power converters having an electric circuit which converts a source of direct current (DC) from one voltage level to another, by storing the input energy temporary and subsequently releasing the energy to the output at a different voltage level.\nThe storage of the input energy may be in either magnetic field storage components (inductors, transformers) or electric field storage components (capacitors). DC-to-DC converters can be designed to transfer power in only one direction, from the input to the output. However many all DC-to-DC converter topologies can be made bi-directional which can move power in either direction which is useful in applications requiring regenerative braking.\nThe amount of power flow between the input and the output of the DC-to-DC converter can be controlled by adjusting the duty cycle (which is identified as a ratio of on/off time of a switch in the DC-to-DC converter). Usually, this is done to control the output voltage, the input current, the output current, or to maintain a constant power.\nIt is desirable to use DC-to-DC converters with minimized switching losses, lower inductor current ripple and low switch voltage ratings, as well as reduced output voltage ripples, and enhanced reliability and autonomity of the operation in plug-in electric vehicles.\nIt is also desirable to provide various configurations of power electronic interfaces between the Energy Storage Sub-Systems and the Propulsion Machines (including electric traction system), in the plug-in electric vehicles, with efficient dynamic splitting of power output from the Energy Storage Sub-Systems to power the Propulsion Machines in accordance with their operational power requirements in a dynamic fashion in order to attain an extended ESS usability and lifetime and to provide the power flow control flexibility."}
-{"text": "1 Field of the Invention\nThis invention relates to an improved system for integrating automatic outbound dialer functions with automatic call distribution functions, and more particularly to a system that balances the number of agents assigned to these functions while maintaining the inbound call waiting time within prescribed limits.\n2 Description of the Prior Art\nMany call centers provide both an automatic inbound telephone call distribution function and an automatic outbound call function. As will be appreciated by those skilled in the art, there have been proposals in the prior art to link inbound automatic call distribution and automatic outbound calling in order to improve staffing efficiency. Inbound call distribution typically has peaks and valleys in its load, since the demand is generated by outside callers. By linking inbound and outbound call functions, outbound agents can be switched to inbound duty during peak inbound demand periods and switched back to outbound duty during slack periods in inbound demand, thus improving overall staffing efficiency.\nIn a typical prior art system, an inbound performance parameter is monitored based upon statistics tracked by the inbound call distributor; for example, number of calls in the inbound queue or average time to answer an inbound call. Target values, and upper and lower thresholds are established for the inbound performance parameter; for example a five second target to respond to incoming calls with an upper threshold of seven seconds and a lower threshold of three seconds. The assumptions of the thresholds is; if upper threshold is exceeded, the performance is unacceptable and more inbound agents are needed; and (b) if the lower threshold is exceeded, the inbound function is considered to be overstaffed and overall efficiency would be improved by transferring inbound agents to outbound operations.\nIn the prior art, so called blend agents are typically transferred, as they become available, when the upper or lower threshold values are exceeded. Agents continue to be transferred until the performance is between the upper and lower threshold limits, when the transfer stops. U.S. Pat. No. 5,425,093, assigned to the assignee of this application and incorporated herein by reference, discloses the concept of providing hysteresis to provide added stabilization to agent transfer. With hysteresis, a higher threshold value is used to trigger the transfer of agents as the upper range of acceptable performance is being exceeded than the threshold value used to stop transfer as the performance moves toward the acceptable range. Similarly for the lower limit; the trigger value to start transfer is lower as measured performance leaves the range than to stop transfer as the performance moves toward the range. However, this and other prior art call center balance algorithms for determining the transfer of agents between inbound and outbound functions operate in response to current inbound performance, and have not proven to be altogether satisfactory in operation."}
-{"text": "In heat assisted magnetic recording, information bits are recorded on a data storage medium at elevated temperatures, and the data bit dimensions can be determined by the dimensions of the heated area in the storage medium or the dimensions of an area of the storage medium that is subjected to a magnetic field. In one approach, a beam of light is condensed to a small optical spot onto the storage medium to heat a portion of the medium and reduce the magnetic coercivity of the heated portion. Data is then written to the reduced coercivity region.\nOne example of a recording head for use in heat assisted magnetic recording generally includes a write pole and a return pole magnetically coupled to each other through a yoke or pedestal, and a waveguide for focusing light onto the storage medium. The waveguide is positioned between the write pole and the return pole. Light from a laser is coupled into the waveguide and transmitted through the waveguide to produce a tightly focused laser spot on the storage media. The power in this spot determines the width of the recorded track and the quality of the recorded transitions.\nIt would be desirable to have a method to measure or monitor the optical power once the head is assembled in a drive."}
-{"text": "This invention relates generally to x-ray imaging processes and systems that use two or more x-ray spectra.\nConventional x-ray imaging uses a single kVp during a scan and represents an observed object in terms of its attenuation. However, this information is not sufficient to precisely characterize the observed object, since attenuation is energy dependent. In the diagnostic energy range, x-ray attenuation is the combination of two photon-matter interactions: the photoelectric effect and Compton scattering. These two interactions are energy-dependent and reflect the effective atomic number and electron density of the object. Accordingly, two measurements at distinct energy spectra can separate the attenuation information into these two basic components, which allows for better identification of the materials present in the object, contrast material specific images, or a weighted sum of the two interactions.\nLarger spectral separation of the low and high energy spectra is favorable to image variance deduction, therefore it contributes to reducing the total radiation dose of each scan while maintaining the same noise level of the output image."}
-{"text": "The invention relates generally to improvements in the lacing assembly of a shoe, and more particularly to an improved lacing assembly for an athletic shoe in which lower and upper vamp sections can be selectively adjusted and maintained to fit different portions of the foot independently while utilizing a single lace for and extending between the lower and upper vamp sections.\nIn applicant's prior U.S. Pat. No. 3,546,796, an athletic shoe is disclosed having separate lower and upper vamp sections. A pair of laces are provided, one lace for each vamp section. While each vamp section can be adjustably fitted and maintained by its own individual lace, it has been found to be disadvantageous to use a pair of separate laces for each shoe. For example, when untied, a pair of laces presents four loose ends that can become tangled and must be stored out into appropriate pairs before each vamp section can be adjustably fitted and tied. Moreover, when a pair of laces are tied, each lace provides a bow, thereby providing a pair of bows on each shoe that can become tangled and which present an unconventional and unacceptable appearance."}
-{"text": "1. Technical Field\nThe present invention relates to an electromagnetic switch. In particular, the present invention relates to an electromagnetic switch provided with a single cylindrical frame and two electromagnets which are accommodated in the frame in parallel with each other in the axial direction of the solenoids.\n2. Related Art\nOne of conventional electromagnetic switches used for a starter of a vehicle is disclosed in Japanese Patent Laid-open Publication No. JP-A-2009-191843.\nThis electromagnetic switch has a first electromagnet for moving the pinion gear of the starter to a ring gear of an engine, and a second electromagnet for closing and opening a main switch. This switch additionally has a cylindrical frame in which the two electromagnets are set in parallel with each other in the axial direction of the switch, and a plastic cover which covers an open end of the frame and is fixed to the frame.\nFIG. 10 shows a schematic view of the electromagnetic switch viewed in the axial direction.\nTwo terminal bolts 110, 120 connected to a motor circuit are fixed to the plastic cover 100. A first exciting terminal 130 to be connected to a first coil of the first electromagnet is shown at the right side of the two bolts, and a second excitation terminal 140 to be connected to a second coil of the second electromagnet is shown at the left side of the two bolts in FIG. 10.\nThe first coil of the electromagnetic switch in JP-A-2009-191843 includes an attracting coil and a holding coil. In addition to the first and second excitation terminal 130, 140, the electromagnetic switch has a negative terminal 150 which is connected to a negative end of the attracting coil. The negative terminal 150 is projected from the plastic cover 100 to the axial direction, electrically connected to the terminal bolt 120 of the motor through an attachment clasp.\nAs shown in FIG. 10, the first and second excitation terminals 130, 140 face each other across the two terminal bolts 110, 120. Therefore it is necessary to wire up one of the excitation terminals 130, 140 before attaching the starter to the engine block. For example, when the engine block and the second excitation terminal 140 are arranged at the left side of the plastic cover 100 in FIG. 10, it is difficult to wire up the second excitation terminal 140 with the starter attached to the engine block. Accordingly, it is required to wire up the second excitation terminal 140 before attaching the starter to the engine block and to wire up the first excitation terminal 130 after attaching the starter to the engine block. This needs twice work of wiring before and after attaching the starter to the engine, and work efficiency is low.\nThe positional relationship of terminals to be wired up is different according to the engine model. It is necessary to change the position of the first and second excitation terminals 130, 140 according to the difference.\nHere a method of change the position of the first and second excitation terminals 130, 140 by rotating about an axis of the frame is discussed.\nThe terminal which is wired up after attaching the starter to the engine is defined as a basis position terminal (first excitation terminal 130 in the case described above). When the first and second excitation terminal 130, 140 are rotated about the axis with respect to the frame for changing the position of the basis position terminal from lower right to lower left, the positions of the first and second excitation terminal 130, 140 move from the position shown in FIG. 11 to the position shown in FIG. 12 with the positional relationship among the negative terminal, the first and second excitation terminal 130, 140 kept.\nIn the above described case, the negative terminal 150 is positioned in a contact space. The contact space is the space where a main switch is arranged. The main switch has a pair of fixed contacts which are connected to the two terminal bolts 110, 120, and a movable contact which closes and opens between the pair of fixed contacts. The electromagnetic switch in JP-A-2009-191843 is configured so that the first and second excitation terminals 130, 140 face each other across the contact space. In such an electromagnetic switch, rotating the first and second excitation terminals 130, 140 about the axis cannot change the position of the basis position terminal to the desired position often.\nIn the case of changing the position of the basis position according to the engine model and the like, for example as shown in FIG. 13 and FIG. 14, it is required to set the appropriate positional relationship between the first and second excitation terminals 130, 140 such that the other terminal wired up in advance (the second excitation terminal in FIG. 12) and the negative terminal 150 don't reach to the contact space. The change needs changing a position defined as a taking position for taking the first coil from a bobbin around which the first coil is wound. Therefore setting a plurality of variations according to different positions of terminals needs a plurality kind of bobbins according to the variations.\nThis cause the number of the component and the equipment used for setting the terminals in the electromagnetic switches according to a plurality of variations to increase. As a result, the production cost increases."}
-{"text": "In a series of previous issued and pending U.S. Patent documents inventors named in connection with the present document and various groupings of their colleagues have disclosed a family of electronic warfare radio receivers that are usable in obtaining military and otherwise useful information from an incoming microwave frequency radio signal and have also disclosed significant components of such receivers. Early in this sequence of patent documents the name \u201cmonobit receiver\u201d was used in referring to both one embodiment of such a receiver and subsequently as a generic or family name for such receivers. Each of these receivers as preferably embodied is of a wide band nature and employs the Fourier transformation as a tool for identifying characteristics of the incoming radio frequency signal. Additionally each of these receivers has the attribute of employing a simplified multiplication arrangement in making use of the Fourier transformation results in order that a simple and low cost realization of the receiver can be achieved.\nIn each of these simplified multiplications the Kernel function portion of the Fourier transformation is represented by a unit value, a magnitude of one or near one, in order to avoid the mechanization complexity of full fledged multiplication algorithm during the receiver step when complex numbers representing an input signal are multiplied by a second set of complex numbers representing the Fourier transformation Kernel function. Generally these simplified multiplications have in the past been accomplished by insuring the second set of complex numbers, i.e., the Fourier transformation results, involve real and imaginary values of unity or near unity. As implied by the name \u201cmonobit receiver\u201d the earliest of these simplified Kernel function receivers employs a Kernel function approximation representable by a single binary bit of information while the later versions of the receiver and indeed the present invention involve Kernel function arrangements requiring a greater number of bits, i.e., three or four bits. In view of these greater numbers of bits it is of course necessary to accomplish an actual multiplication operation involving the Fourier transformation results however for feasibility purposes some form of simplified and readily accomplished multiplication is needed.\nIn order to better appreciate the present invention it may be helpful to regress briefly and consider the evolution of simplified Kernel function realizations in our work. Such evolution commences with the realization of a Fourier transformation using only multiplication by unity or in essence no multiplication in the Fourier transformation computation algorithm. Kernel function realization in this manner is first disclosed in a U.S. patent of Tsui et al., a patent numbered U.S. Pat. No. 5,917,737, a patent wherein Kernel function values are located on a circle of unit radius at angular locations of 5/4, 3 5/4, 5 5/4 and 7 5/4 radians i.e., at locations displaced by forty five degrees from the real and imaginary axes of a coordinate axis plot. Actual coordinate axis locations of 1+j, 1\u2212j, \u22121+j and \u22121+j are used for the Kernel functions disclosed in the U.S. Pat. No. 5,917,737 patent. For reference simplification purposes this and the several other patents issued to various combinations of persons from our laboratory are herein referred-to as \u201cour patent\u201d notwithstanding the differing group of inventors, with perhaps one common inventor, associated with each patent.\nOur U.S. Pat. No. 5,793,323 therefore relates to the U.S. Pat. No. 5,917,737 patent in that it discloses a single integrated circuit chip arrangement for a monobit receiver employing the approximated Kernel function of the U.S. Pat. No. 5,917,737 patent. This patent also discloses several implementation compromises possible in applying the simplified Kernel function to receiver apparatus and provides a receiver that can cover a one gigahertz spectrum with a frequency resolution of 9.77 megahertz.\nOur U.S. patent application Ser. No. 09/917,589 filed on Jul. 30, 2001 also relates to the U.S. Pat. No. 5,917,737 patent in that it discloses the use of a monobit receiver in combination with several bandpass filters in order to increase the number of simultaneously processable signals and enhance the dynamic range capability of the overall system.\nThe invention of our U.S. patent application Ser. No. 09/944,616, filed on Sep. 4, 2001, provides a straightforward approach to the enhancement of dynamic range in a monobit family receiver by increasing the number of Kernel function locations used in the Kernel function approximation from the four locations of the U.S. Pat. No. 5,917,737 patent and the U.S. Pat. No. 5,963,164 patent to eight locations. In the application of Ser. No. 09/944,616 Kernel function values located at the 5/4, 3 5/4, 5 5/4 and 7 5/4 radian locations are added to the Kernel function values at 0, 5/2, 5, and 3 5/2 radians with the added four values being slightly increased in magnitude from true unit circle values and in fact having a magnitude of (2)1/2 or 1.414. Some of these eight location Kernel function values depart slightly from unity magnitude however the overall results of the eight-location Kernel function appear favorable.\nAdditional of our U.S. Patent and Trademark Office documents involving Kernel function realizations include the application of Ser. No. 10/115,819, filed on Apr. 3, 2002, and now issued as Statutory Invention Registration H2109. This document is titled \u201cPASSIVE MICROWAVE DIRECTION FINDING WITH MONOBIT FOURIER TRANSFORMATION RECEIVER AND MATRIX COUPLED ANTENNA\u201d\nIn addition to these previous patent documents wherein use of a Kernel function realization of unity or near unity value is employed in order to achieve a simplified multiplication operation the U.S. patent application Ser. No. 10/354,267, filed on Jan. 31, 2003, discloses an electronic warfare receiver wherein a Kernel function realization of other than unity magnitude is used in combination with a simplified multiplication arrangement accomplished with a bit shifting multiplication algorithm. In the of Ser. No. 10/354,267, document the simplified bit shifting algorithm multiplication can be implemented with such minimal complexity and cost as to make the non-unity magnitude of the Kernel function of little consequence. The present invention may be considered as an extension of this simplification multiplication through shifting concept.\nThe present invention may also start with the initial realization that a significant difficulty with the monobit electronic warfare receiver results form the fact that the two signal instantaneous dynamic range achieved in such a receiver is undesirably low. The instantaneous dynamic range relates to the receiver's capability to detect two simultaneous signals of different amplitude and can be an important consideration in a receiver use situation especially when a hostile adversary is involved. Because of the non-linear property of the approximated Kernel function in the Tsui et al. U.S. Pat. No. 5,917,737 patent receiver, the two signal instantaneous dynamic range of this receiver is limited to about 2 dB. Using a different Kernel function in the later tribit receiver of the Ser. No. 09/944,616 application can improve this instantaneous dynamic range to about 10 dB. A receiver arrangement according to the present invention can further improve this dynamic range characteristic. Each of the patent documents, issued or pending, identified in this BACKGROUND OF THE INVENTION DISCUSSION is hereby incorporated by reference herein.\nThe present patent document continues in this line of Kernel function and Kernel function receiver patents and thus provides additional insight into the simplified Kernel function utilization art. In the present invention additional improvement of the dynamic range achieved in the monobit family of receivers is pursued while relaxing the previously perceived need to limit multiplications to unity magnitude values."}
-{"text": "It has long been known in the art that a contact lens must have surfaces that have a certain degree of hydrophilicity in order to be wet by tears thus providing unblurred vision.\nSoft, hydrophilic contact lenses, in addition to being wettable, provide comfort to the wearer but lack the ability to correct certain visual deficiencies such as astigmatism since they tend to conform to the shape of the corneal surface.\nOften hydrophilic monomers can be added to a mixture of comonomers in the formation of contact lenses so that upon polymerization optically clear contact lenses result which have a certain degree of hydrophilicity. As the hydrophilic monomer content increases where it is added directly to the lens composition, the physical characteristics of the lenses are affected by the increased hydration propensity of the polymeric composition.\nIn some cases, it has been known to treat a formed contact lens with a polymerizable hydrophilic monomer to form a surface coating of hydrophilic polymer grafted to an otherwise hydrophobic polymer surface. Although effective, this method of increasing the hydrophilic character of the lens surface can suffer from involved and difficult manufacturing procedures as well as lack of permanence.\nPresent rigid and soft contact lenses sometimes retain water on their surfaces through secondary chemical bonding and as a consequence only a very thin layer of water molecules is present between the eye and the contact lens.\nSoft lenses are inherently comfortable but often times, as with hard lenses, suffer from brief surface dryness between eye blinks. State of the art technology teaches that a water soluble neutral polymer may be applied to the surfaces of a hard contact lens to provide a \"cushion\" layer between the lens and the eye which is equated with increased wettability as well as wearer comfort and tolerance.\nDissipation of the \"cushion\" layer occurs rapidly in most prior art constructions, since there is little specific interaction between the mobile polymer in this layer and the lens surface. As a result the wearer begins to feel discomfort and must rewet the lens surfaces."}
-{"text": "1. Field of Technology\nThe present disclosure relates to equipment and techniques for melting metals and metallic alloys (hereinafter \u201calloys\u201d). The present disclosure more specifically relates to equipment and techniques utilizing electrons to melt metals and alloys and/or the materials while in a molten state.\n2. Description of the Background of the Technology\nThe alloy melting process involves preparing a charge of suitable materials and then melting the charge. The molten charge or \u201cmelt\u201d may then be refined and/or treated to modify melt chemistry, remove undesirable components from the melt, and/or affect the microstructure of articles cast from the melt. Melting furnaces are powered by either electricity or the combustion of fossil fuels, and selection of a suitable apparatus is largely influenced by the relative costs and applicable environmental regulations, as well as by the identity of the material being prepared. A variety of melting techniques and apparatus are available today. General classes of melting techniques include, for example, induction melting (including vacuum induction melting), arc melting (including vacuum arc skull melting), crucible melting, and electron beam melting.\nElectron beam melting typically involves utilizing thermo-ionic electron beam guns to generate high energy substantially linear streams of electrons which are used to heat the target materials. Thermo-ionic electron beam guns operate by passing current to a filament, thereby heating the filament to high temperature and \u201cboiling\u201d electrons away from the filament. The electrons generated from the filament are then focused and accelerated toward the target in the form of a very narrow, substantially linear electron beam. A type of ion plasma electron beam gun also has been used for preparing alloy melts. Specifically, a \u201cglow discharge\u201d electron beam gun described in V. A. Chernov, \u201cPowerful High-Voltage Glow Discharge Electron Gun and Power Unit on Its Base\u201d, 1994 Intern. Conf. on Electron Beam Melting (Reno, Nev.), pp. 259-267, has been incorporated in certain melting furnaces available from Antares, Kiev, Ukraine. Such devices operate by producing a cold plasma including cations which bombard a cathode and produce electrons that are focused to form a substantially linear electron beam.\nThe substantially linear electron beams produced by the foregoing types of electron beam guns are directed into the evacuated melting chamber of an electron beam melting furnace and impinged on the materials to be melted and/or maintained in a molten state. The conduction of electrons through the electrically conductive materials quickly heats them to a temperature in excess of the particular melting temperature. Given the high energy of the substantially linear electron beams, which can be, for example, about 100 kW/cm2, linear electron beam guns are very high temperature heat sources and are readily able to exceed the melting and, in some cases, the vaporization temperatures of the materials on which the substantially linear beams impinge. Using magnetic deflection or similar directional means, the substantially linear electron beams are rastered at high frequency across the target materials within the melting chamber, allowing the beam to be directed across a wide area and across targets having multiple and complex shapes.\nBecause electron beam melting is a surface heating method, it typically produces only a shallow molten pool, which can be advantageous in terms of limiting porosity and segregation in the cast ingot. Because the superheated metal pool produced by the electron beam is disposed within the high vacuum environment of the furnace melting chamber, the technique also beneficially tends to degas the molten material. Also, undesirable metallic and non-metallic constituents within the alloy having relatively high vapor pressures can be selectively evaporated in the melting chamber, thereby improving alloy purity. On the other hand, one must account for the evaporation of desirable constituents produced by the highly-focused substantially linear electron beam. Undesirable evaporation must be factored into production and can significantly complicate alloy production when using electron beam melting furnaces.\nVarious melting and refining methods involve the electron beam melting of feed stocks using thermo-ionic electron guns. Drip melting is a classic method used in thermo-ionic electron beam gun melting furnace for processing refractory metals such as, for example, tantalum and niobium. Raw material in the form of a bar is typically fed into the furnace chamber and a linear electron beam focused on the bar drip-melts the material directly into a static or withdrawal mold. When casting in a withdrawal mold, the liquid pool level is maintained on the top of the growing ingot by withdrawing the ingot bottom. The feed material is refined as a result of the degassing and selective evaporation phenomena described above.\nThe electron beam cold hearth melting technique is commonly used in the processing and recycling of reactive metals and alloys. The feedstock is drip melted by impinging a substantially linear electron beam on an end of a feedstock bar. The melted feedstock drips into an end region of a water-cooled copper hearth, forming a protective skull. As the molten material collects in the hearth, it overflows and falls by gravity into a withdrawal mold or other casting device. During the molten material's dwell time within the hearth, substantially linear electron beams are quickly rastered across the surface of the material, retaining it in a molten form. This also has the effects of degassing and refining the molten material through evaporation of high vapor pressure components. The hearth also may be sized to promote gravity separation between high-density and low-density solid inclusions, in which case oxide and other relatively low-density inclusions remain in the molten metal for a time sufficient to allow dissolution while high density particles sink to the bottom and become trapped in the skull.\nGiven the various benefits of conventional electron beam melting techniques, it would be advantageous to further improve this technology."}
-{"text": "1. Field of the Invention\nThis invention relates to helmet mounted displays. In this specification the term \"helmet\" means any article worn by an observer on his/her head and includes protective helmets such as those worn by aircraft pilots or vehicle drivers and simple instrument supporting head bands such as those worn by surgeons and other medical workers. The term \"helmet mounted display\" shall mean apparatus attached to, located on, or forming part of such a helmet and for presenting images or information to the wearer of the helmet.\n2. Description of the Related Art\nFor combat aircraft use, flight testing and medical examination have confirmed that binocular display to the pilot both aids performance and also reduces brain fatigue.\nAt present, binocular presentation is provided by either having:\n(i) two sources of display or PA1 (ii) one source with optical train containing a split prism.\nThe adoption of solution (i) doubles the number of devices on the helmet providing the display and hence increases the helmet mass.\nThe adoption of solution (ii) halves the level of brightness of the display as inherently the split prism has a 50% transmission characteristic.\nConsequently both the existing techniques have operational disadvantages. Firstly, excess mass and the adverse physiological effect under g forces; and secondly, low brightness and hence loss of display under high ambient light conditions.\nThere is alternatively often a requirement to provide the pilot with information f rom two similar or dissimilar data sources simultaneously. Such sources may be for example cathode ray tube, dot matrix, or conventional optical displays. Currently the pilot may view one source only at a time."}
-{"text": "Embodiments of the invention relate to neuromorphic and synaptronic computation and in particular, mapping neural dynamics of a neural model on to a lookup table.\nNeuromorphic and synaptronic computation, also referred to as artificial neural networks, are computational systems that permit electronic systems to essentially function in a manner analogous to that of biological brains. Neuromorphic and synaptronic computation do not generally utilize the traditional digital model of manipulating 0s and 1s. Instead, neuromorphic and synaptronic computation create connections between processing elements that are roughly functionally equivalent to neurons of a biological brain. Neuromorphic and synaptronic computation may comprise various electronic circuits that are modeled on biological neurons.\nIn biological systems, the point of contact between an axon of a neuron and a dendrite on another neuron is called a synapse, and with respect to the synapse, the two neurons are respectively called pre-synaptic and post-synaptic. The essence of our individual experiences is stored in conductance of the synapses. The synaptic conductance changes with time as a function of the relative spike times of pre-synaptic and post-synaptic neurons, as per spike-timing dependent plasticity (STDP). The STDP rule increases the conductance of a synapse if its post-synaptic neuron fires after its pre-synaptic neuron fires, and decreases the conductance of a synapse if the order of the two firings is reversed."}
-{"text": "1. The Field of the Invention\nThe present invention relates to sodium hypochlorite and other hypohalite or hypohalous acid compositions (e.g., aqueous solutions) that include a cationic polymer, as well as methods of making and using such compositions.\n2. Description of Related Art\nConsumers recognize sodium hypochlorite as a highly effective cleaning, bleaching and sanitizing agent that can be widely used in cleaning and sanitizing various hard and soft surfaces, in laundry care, etc. Often, sodium hypochlorite solutions in a relatively concentrated form are delivered to the consumer, who then dilutes them. A wide variety of automated washing machines which can control the addition of sodium hypochlorite to laundry washwater are also available. When used in washing machines, the sodium hypochlorite solution will typically come in contact with the detergent being used in the laundering process. Complex interactions between sodium hypochlorite and components of the detergent, such as surfactants, builders, enzymes, fragrances, dyes, fabric brighteners, fabric softener components, and anti-redeposition polymers and fabric care polymers can occur, and are thought to result in some cases in a reduction in the stain removal and bleaching performance of sodium hypochlorite, or reduction in the aesthetics or detergency performance of the detergent formulation used. Such complex interactions are commonly observed when the stain removal and whitening performance of sodium hypochlorite solutions is measured in the presence of different detergent formulations, whereby some detergent formulations are more compatible with the sodium hypochlorite solution than others. For example, fabric brighteners are often less effective or ineffective when contacted with the bleaching agent of the bleach solution.\nIn the case of manual washing of laundry, where sodium hypochlorite is employed, the relative volume of water to fabrics may be reduced, and hence the levels of suspended soils, both particulate and oily, may be significantly higher than in an automatic washing machine.\nSodium hypochlorite solutions designed to be diluted in such laundering processes may contain a variety of adjuvants to enhance the stability of the hypochlorite or enhance the aesthetics of use of the product (e.g., lower odor). Fragrances, hydrotropic materials (e.g., such as sodium xylene sulfonates), buffers, certain electrolytes (e.g., alkali metal halides, etc.) may be added to hypohalite solutions for such purposes.\nHowever, variation in the overall consumer-perceivable fabric whitening and stain removal still exists when hypohalite solutions are used in combination with commercial detergent formulations. Thus, there continues to be a need for liquid compositions including hypohalite species, which compositions would exhibit improved or more consistent cleaning, bleaching, and whitening performance when used with various modern detergent formulations, or when used under conditions of high loads of suspended soils."}
-{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent it is described in this background section, as well as aspects of the description that may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure.\nIn an embedded system, there are typically constraints on hardware resources thereby limiting processing capabilities. This limits the speed in which the system is able to load all the necessary software components to become useable. In particular, embedded systems with a graphical user interface (GUI) will take an inordinate amount of time as GUIs are typically resource intensive.\nTherefore, an unaddressed need exists in the art to address the aforementioned deficiencies and inadequacies."}
-{"text": "This invention relates to novel polymer product compositions, polymerization feed compositions and monomer purification therefor. In particular, it relates to a crosslinked, high modulus, high impact strength, thermoset polymer of polycycloolefin units which is formed via a metathesis-catalyst system.\nDistillation is commonly used in the preparation of polycycloolefins, for example dicyclopentadiene. Various purities of dicyclopentadiene are available. The invention preferably uses economically available purified dicyclopentadiene which is readily polymerized. Purified dicyclopentadiene forms a substantially cross-linked thermoset polymer when polymerized as disclosed herein.\nAny good thermoset polymer should meet at least two criteria. It should have desirable physical properties and it should lend itself to easy synthesis and forming. Among the most desirable physical properties for many polymers is a combination of high impact strength and high modulus. A standard test for impact strength is the notched Izod impact test, ASTM No. D-256. For an unreinforced thermoset polymer to have good impact strength, its notched Izod impact should be at least 1.5 ft. lb./in. notch. It is desirable that this good impact strength be combined with a modulus of at least about 150,000 psi at ambient temperature. Thermoset polymers with high impact strength and high modulus find useful applications as engineering plastics in such articles of manufacture as automobiles, appliances and sports equipment. Among the critical factors in the synthesis and forming of a thermoset polymer are the conditions required to make the polymer set up or gel. Many thermoset polymers require considerable time, elevated temperature and pressure, or additional steps after the reactants are mixed before the setting is complete.\nA thermoset homopolymer having high impact strength and high modulus has been described by Klosiewicz in U.S. Pat. Nos. 4,400,340, 4,469,809 and 4,436,858 (with plasticizer) and by Leach in 4,458,037 (a foam). Characteristics of thermoset polymers include insolubility in common solvents such as gasoline, naphtha, chlorinated hydrocarbons, and aromatics as well as resistance to flow at elevated temperatures.\nWork has been done on the metathesis copolymerization of dicyclopentadiene with one or more other monomers to produce soluble copolymers. This copolymer formation has resulted in the production of unwanted insoluble by-products. U.S. Pat. No. 4,002,815, for instance, which teaches the copolymerization of cyclopentene with dicyclopentadiene, describes an insoluble by-product and suggests that the by-product could be a gel of a dicyclopentadiene homopolymer.\nSome other work, usually in an attempt to produce soluble polymers, has been done on the metathesis polymerization of dicyclopentadiene. Japanese unexamined published patent applications KOKAI Nos. 53-92000 and 53-111399 disclose soluble polymers of dicyclopentadiene. Several syntheses of soluble polymers of dicyclopentadiene have produced insoluble by-products. Takata et al, J. Chem. Soc. Japan Inc. Chem. Sect., 69, 711 (1966), discloses the production of an insoluble polymerized dicyclopentadiene by-product from the Ziegler-Natta catalyzed polymerization of dicyclopentadiene; Oshika et al, Bulletin of the Chemical Society of Japan, discloses the production of an insoluble polymer when dicyclopentadiene is polymerized with WCl.sub.6, AlEt.sub.3 /TiCl.sub.4 or AlEt.sub.3 /MoCO.sub.5 ; and Dall Asta et al, Die Makromolecular Chemie 130, 153 (1969), discloses an insoluble by-product produced when a WCl.sub.6 /AlEt.sub.2 Cl catalyst system is used to form polymerized dicyclopentadiene.\nIn U.S. Pat. No. 3,627,739 ('739), dicyclopentadiene is gelled with unactivated catalyst and then heated for an hour.\nPampus et al in U.S. Pat. No. 3,873,644 discloses a method of producing graft polymers from a cyclic olefin to obtain thermoplastic products of high impact strength.\nOfstead in U.S. Pat. Nos. 4,020,254 and 3,935,179 disclose metathesis polymerization of cycloolefins to obtain a rubbery polymer. Streck et al in U.S. Pat. Nos. 3,974,092 and 3,974,094 discloses a catalyst for preparation of polyalkenamers. Each discloses polyalkenamers of low reduced melt viscosity. Percent gel obtained using the system of Streck et al is very low. A typical percent gel is about 2 to 6% with 14 being the highest value disclosed. Wilkes in U.S. Pat. No. 3,084,147 discloses thermalpolymerization of dicyclopentadiene. The thermal polymerization is carried out in a nonpolymerizable solvent at about 500.degree. to 550.degree. F.\nStafford in U.S. Pat. No. 3,446,785 discloses a polymerization of olefins. The polymerization discloses the production of a product which is a viscous liquid polymer or a brittle, semi-solid polymer.\nNutzel et al in U.S. Pat. No. 3,684,787 discloses preparation of polyalkenamers. The polymers can be isolated by pouring the solution of polymer into 3 to 5 times its quantity of a solution of a lower alcohol in which an age resister is dissolved. Alternatively, the solution can be introduced into boiling water and the solvent removed with steam. Mensel states in column 1, line 19 that crosslinked products are of no industrial interest. The polymers obtained have a rubber-like character.\nVergne et al in U.S. Pat. No. 3,557,072 discloses plastomers derived from dimethanooctahydronaphthalene and their method of manufacture. Amorphous polymers are obtained in the form of a hard white mass which is washed with methanol and then ground and dried to a white powder.\nTenney et al and Tenney alone in U.S. Pat. Nos. 4,136,247, 4,136,248, 4,136,249 and 4,178,424 disclose ring-opened cycloolefin polymers and copolymers. These polymer products are thermoplastics which can be thermoformed.\nMinchak and Minchak et al respectively in U.S. Pat. Nos. 4,002,815 and 4,380,617 each disclose polymerization of cycloolefins to form polymers which may be isolated by precipitation using an alcohol or by steam or hot water stripping. The polymers produced have inherent viscosities from about 0.1 to about 10 and are greater than 90% soluble in solvent. Viscous cements and plastic polymers solidify in from about 30 minutes to 180 minutes. The reaction is short stopped in less than 2 hours by addition of an alcohol. Minchak in U.S. Pat. No. 4,426,502 discloses bulk polymerization of cycloolefins by reaction injection molding in less than about 2 minutes using an organoammonium molybdate or tungstate catalyst.\nDeWitt et al in U.S. Pat. No. 4,418,179 discloses impact modification cycloolefins by polymerization using an organoammonium molybdate or tungstate catalysts in the less than 2 minutes using reaction injection molding.\nOshika et al in the Bulletin of the Chemical Society of Japan, line 41, pages 211-217 (1968) discloses ring opening polymerization of norbornene and its derivatives by MoCl.sub.5, WCl.sub.6 and ReCl.sub.5 catalysts. Dark crude polymer is obtained which is dissolved and reprecipitated with methanol.\nU.S. Pat. No. 4,002,815 discloses the use of a metathesiscatalyst system which employs a dialkylaluminum iodide, an alkylaluminum diiodide or a mixture of trialkylaluminum compounds with elemental iodine to produce substantially gel-free copolymers of cyclopentene and dicyclopentadiene.\nU.S. Pat. No. 4,069,376 discloses the use of a three component catalyst comprised of a soluble tungsten compound, a dialkylaluminum chloride or alkylaluminum dichloride, and a dialkylaluminum iodide or alkylaluminum diiodide to produce substantially gel-free norbornene-dicyclopentadiene copolymers.\nU.S. application Ser. No. 526,835 filed Aug. 26, 1983 and assigned to the same assignee, discloses a cellular crosslinked poly(dicyclopentadiene) which is made with a metathesis-catalyst system. The cellular polymer is made by injecting the catalyst system, which includes an alkylaluminum activator, into a reaction vessel which is preheated, preferably to a temperature from about 100.degree. C. to about 125.degree. C.\nNot only is it desirable that the thermoset polymer have high impact strength, but it is also desirable that it be easily synthesized and formed. A reaction injection molding (sometimes hereinafter referred to as RIM) process achieves this second goal by in-mold polymerization. The process involves the mixing of two or more low viscosity reactive streams. The combined streams are then injected into a mold where they quickly set up into a solid infusible mass. For a RIM system to be of use with a particular polymer, certain requirements must be met: (1) the individual streams must be stable and must have a reasonable shelf-life under ambient conditions; (2) it must be possible to mix the streams thoroughly without their setting up in the mixing head; (3) when injected into the mold, the materials must set up to a solid system rapidly; and (4) any additives-fillers, stabilizers, pigments, etc. - must be added before the material sets up. Therefore, the additives selected must not interfere with the polymerization reaction.\nInterpenetrating polymer networks (IPNs) are a type of polymer alloy consisting of two (or more) crosslinked polymers. They are more-or-less intimate mixtures of two or more distinct, crosslinked polymer networks held together by permanent entanglements with few, if any covalent bonds between the polymers. The entanglements in IPNs must be of a permanent nature and are made so by self-crosslinking of the two polymers. They are introduced either by swelling a crosslinked polymer with monomer and the crosslinking agent of another polymer, and curing the swollen polymer in situ or by mixing the linear polymer, prepolymers, or monomers in some liquid form solution, or bulk, together with crosslinking agents, evaporating the vehicle (if any), and curing the component polymers simultaneously.\nIPNs possess several interesting characteristics in comparison to normal polyblends. Formation of IPNs is the only way to intimately combine crosslinked polymers, the resulting mixture exhibiting (at worst) only limited phase separation. Normal blending or mixing of polymers results in a multiphase morphology due to the well-known thermodynamic incompatibility of polymers. However, if mixing is accomplished simultaneously with crosslinking, phase separation may be kinetically controlled by permanent interlocking of entangled chains.\nMolecular and supramolecular perspectives are believed to be important in considering interpenetrating networks. On linked polymer networks intermeshed through molecular chain segment entanglements. These are totally non-separable without chain breaking (i.e. degradation). A semi-interpenetrating network (SIPN) then is a crosslinked network having a non-crosslinked polymer dispersed through it on a molecular scale. These in theory are separable without bond breaking. Neither the interpenetrating network nor the SIPN represents a heterophase polymer structure in a morphological or supramolecular sense.\nThe supramolecular scale microscopically observable heterophases are present SIPN and IPN. In this context an interpenetrating network is any system of two polymers where the supramolecular structure involves two intermeshed continuous polymer phases irrespective of whether or not they are crosslinked.\nBecause molecular mixing of dissimilar uncrosslinked polymers is thermodynamically unfavorable (for entropic reasons) it is unusual to form a stable, molecular IPN from two uncrosslinked polymers. They generally will separate into a heterophase polymer system which might or might not be a supramolecular IPN.\nThe forming of supramolecular interpenetrating networks of two uncrosslinked polymers is called polymer blends. \"Co-reacted\" polymers is a process which results in block copolymers and sometimes results in heterophase polymer systems which could exhibit supramolecular IPN morphology, depending on block length."}
-{"text": "The use of various steroids to increase the intensity of the heartbeat is known. A number of steroid glycosides have been found useful in cardiotherapy, but the corresponding aglycons (genins), which share the characteristic steroid structure of the glycosides are convulsive poisons.\nFor purposes of definition, a cardiotonic glycoside is one which increases the intensity of the heart beat but may decrease the rate of heartbeat. These glycosides may be designated by the following general formula: ##SPC1##\nThe corresponding aglycons have a hydroxy group at position 3 rather than a glycoside conjugate.\nA wide number of sugars may be employed to form cardiotonic glycosides, as long as the steroid nucleus is a cardenolide which contains a five membered lactone ring at the C.sub.17 position. Particular cardiotonic glycosides useful in the treating of heart disceases and which are readily available include digoxin, digitoxin and gitoxin. These compounds differ at the C.sub.12 and C.sub.16 position of the steroid (R.sub.1 and R.sub.2), in that digitoxin contains a hydrogen substituent at each position, digoxin contains a hydroxyl at C.sub.12 and a hydrogen at C.sub.16 while gitoxin contains a hydrogen at C.sub.12 and hydroxyl at C.sub.16. The glycoside portion of each of these compounds comprises three digitoxose sugars. The 3', 3\", 3'\" and 4'\" positions bear hydroxyl groups while the 5', 5\" and 5'\" bear methyl substituents. Other sugars which may be found in cardiotonic glycosides include the lanatoside series which differs from the digitoxose in that the 3'\" position bears an acetyl while the 4'\" position contains an additional glucose sugar. When this glycoside is conjugated to digoxigenin it is designated lanatoside A, while if it is conjugated with gitoxigenin it is designated lanatoside B, and if it is conjugated with digitoxigenin it is designated lanatoside C.\nOther sugar units which may give a compound which has cardiotonic activity, when they are conjugated to a cardenolide steroid, include rhamnose, antiarose, digitalose, thevetose, talomethylose, cymarose, oleandrose, sarmentose, boivinose and diginose. By combining three of these units with an aglycon, various cardiotonic glycosides may be formed. The aglycons alone are usually convulsive poisons not useful in medicine. For a more thorough discussion of cardiotonic glycosides and aglycons, see \"Rodd's Chemistry of Carbon Compounds\", Second Edition, Vol. 2, part D, chapter 17, and \"Chemistry and Metabolism of Digitalis\", Digitalis, edited by Charles Fisch and Borys Furawicz, Grune & Stratton, Inc., New York, 1969.\nAs previously mentioned, although the glycosides have been found to exhibit much less toxicity than the corresponding aglycon, in many instances the difference between a therapeutic and a toxic dose of a glycoside is only in the order of a few micrograms per kilogram of blood. Digitalis intoxication occurs frequently due to the small differences between therapeutic and toxic quantities. The therapeutic quantity of any cardiotonic glycoside cannot be predicted in advance of medication since each patient may react differently to the dosage, and retain different quantities of the drug in the blood stream.\nDigitalis intoxication occurs with increasing frequency in elderly heart patients, since diseased renal function in the aged can result in higher and sometimes toxic levels of cardiotonic glycoside residue in the blood. Also, there may be some overlapping between therapeutic and toxic blood levels due to susceptability of the myocardium to digitalis. For example, the patient whose myocardium is diseased may show signs of toxicity despite therapeutic blood levels. Other clinical conditions contributing to digitalis intoxication are coronary ischemia and potassium depletion. The problems of digitalis intoxication are compounded in that its symptoms may often be mistaken for fatigue or restlessness.\nGenerally, accepted therapeutic levels for digoxin range from 0.8 to 2.4 ng/ml while toxic levels range from 2.1 to 8.7 ng/ml Toxic versus non-toxic levels of digitoxin also overlap, i.e., 3.0 to 39.0 ng/ml may be therapeutic while 26.0 to 43.0 ng/ml may be toxic. Higher levels of digoxin and digitoxin are generally lethal.\nRadioimmunoassay for cardiotonic glycosides has proved extremely effective in detecting undermedicated, adequately medicated and over-medicated or toxic levels of the drug in blood serum. That sampling technique comprises the use of a radioactive steroid glycoside or aglycon and an antibody which will bind both the radioactive steroid and the cardiotonic glycoside in the blood on a competitive basis. By counting the amount of radioactive material reacting with the antibody, the amount of cardiotonic glycoside in the blood may be calculated. From the results, it can be quickly determined whether the patient has been undermedicated, properly medicated or over-medicated and is suffering digoxin toxicity.\nFor radio labelling of steroids, 125-Iodine or 131-Iodine are preferred, since they exhibit high specific activity, while .sup.14 C and .sup.3 H are also available. The use of compounds exhibiting a high specific activity permits the use of considerably less material for the same counting efficiency. Counting with such isotopes can be achieved by liquid scintillation or gamma ray spectroscopy counting procedures.\nThe sampling procedure preferably consists of the use of an antibody polymer in tablet form, which antibody has been raised in an animal host by the administration of an appropriate antigen. The radioactive steroid and the cardiotonic glycoside isolated in the blood serum then complete with each other for a limited number of binding sites on the antibody to form an insoluble complex. The antigen/antibody complex is then separated via centrifuge techniques and the radioactive portion is measured, as for example, by gamma counter. From the difference in radioactivity in the test specimen, compared to standards of known cardiotonic glycoside concentrations in blood, the level of cardiotonic glycoside in the test sample may be accurately computed.\nFor a more complete discussion of radioimmunoassay principles, see A. R. Midgday, Jr. and G. D. Niswender, in \"Karolinska Symposia on Research Methods in Endocrinology, 2nd Symposium: Steroid Assay by Protein Binding,\" The Reproductive Endocrinology Research Unit, Stockholm (1970), 320-333 and W. D. Odell and W. H. Daughadey, \"Principles of Competitive Protein-Binding Assays,\" J. B. Lippincott Company, Philadelphia (1971). For a complete discussion on the preparation of protein-digoxin conjugates necessary for the production of antibodies, see \"Digoxin-Specific Antibodies\", Butler et al., Proc. Natl. Acad. Sci. U.S. 57:71.\nThe method most commonly used in conjugating a radioactive isotope, such as 125-Iodine, to a steroid is via an electrophilic substitution reaction. Because of the high susceptability of phenol derivatives to electrophilic substitution, since the hydroxyl is an activating group, radioactive steroids have been synthesized by conjugating tyrosine or tyrosine methyl ester (TME) onto the steroid, after which the compound may be reacted with the radioactive isotope (.sup.125 I).\nOliver et al., Journal of Clinical Investigation, 47, 1035 (1968) teach the synthesis of a digitoxigenin-3-0-succinyl-iodine-125 tyrosine methyl ester for use in cardiotonic glycoside radioimmunoassay procedures. Such a compound, although suitable for measuring the level of the drug in blood serum, suffers from the disadvantage that it is extremely difficult and expensive to produce. The succinylation step of the C.sub.3 -hydroxyl group of the genin, which is used at the binding site for the TME and iodine-125, proceeds very slowly and may require reaction time on the order of months. It is believed that the slow reaction time associated with the succinylation of the C.sub.3 -hydroxyl, is due to the fact the C.sub.3 -hydroxyl group is fixed in an axial position and offers a large amount of steric hindrance. Additionally, although such a reaction procedure is suitable for the synthesis of digitoxigenin derivatives, where the C.sub.12 substituent is a hydrogen, it is not per se suitable for the synthesis of radioactive digoxigenin derivatives. Digoxigenin has at the C.sub.12 position and equatorial hydroxyl substituent which will more readily be succinylated than the axial-3-hydroxyl group, since it offers less steric hindrance. Thus succinylation of the 3-hydroxyl of the genin not only proceeds slowly but for many compounds produces little or no yield of a useful immunoreactive derivative unless other positions are blocked. In addition, Stall et al., \"The Specificity of the Digoxin Radioimmunoassay Procedure\", Res. Commun. Chem. Pathol Pharmacol, 4, 503 (1972) shows that digoxigenin has a lower affinity for antibody than digoxin itself. They show a crossreactivity of about 34% wherein the optimum ratio of binding ability of radiolabelled to unlabelled antigen should be 1:1. Thus it would be expected that a radioactive glycoside will bind on a more competitive basis with an antibody than the corresponding aglycon.\nAlthough Draws et al., \"Faster and Easier Radioimmunoassay of Digoxin,\" Clinical Chemistry, Vol. 20, No. 3, 343-347 (1974) discusses the use of a .sup.125 I-tyrosine methyl ester of digoxin, there is no indication that such a compound should be a sugar hydroxyl substituted material or any indication of how such a compound might be synthesized.\nIt is an object of this invention to provide a process of synthesizing a radioactive steroid which will complete on a predictable basis for sites on an antibody with cardiotonic glycosides to be measured.\nIt is another object of this invention to provide a process of synthesizing such radioimmunoassay compounds via an improved and simplified procedure.\nIt is yet another object of this invention to provide a radioimmunoassay compound having a high specific activity and improved stability upon storage."}
-{"text": "The present invention relates to an analytical method for assessing the presence of alkaline sphingomyelinase in the stools or biological fluids of patients in need of such an assessment. The invention also relates to a kit for carrying out the analytical method.\nMore particularly the method of the present invention is an in vitro fluorometric method for detecting alkaline sphingomyelinase which, as will be described in detail hereinbelow, is a marker of serious pathological states such as colon cancer and familial adenomatous polyposis.\nThe enzyme sphingomyelinase (sphingomyelin phosphodiesterase, SMase) catalyzes the hydrolysis of sphingomyelin to ceramide and choline phosphate.\nThree different types of SMase (acidic, neutral and alkaline) have been identified to-date, which occur as several isoforms, as follows:\nlysosomal acidic SMase (A-SMase);\ncytosolic Zn2+-dependent acidic SMase;\nmembrane neutral magnesium-dependent SMase (N\u2014SMase);\ncytosolic magnesium-independent N\u2014SMase; and\nalkaline SMase.\nSMases have been shown to play a role in a wide variety of physiologic and pathological processes, including: lysosomal hydrolysis of endocytosed SM, ceramide mediated cell signaling, atherogenesis, terminal differentiation, cell cycles arrest, apoptosis, inflammation, and the regulation of eukaryotic stress responses.\nIn contrast to acidic and neutral SMase, which are currently present in cells as lysosomal and membrane-bound enzymes, respectively, alkaline SMase exhibits tissue and species difference. In human beings, the alkaline SMase is found in intestinal mucosa and bile. Alkaline SMase starts to appear in the duodenum, reaches a high level in the intestine, especially in the distal part of the jejunum, and occurs in considerable amounts in the colon and rectum. This SMase presents optimal alkaline pH at 9.0, is Mg2+-independent, bile salt-dependent and trypsin-resistant.\nThe pathological importance of alkaline SMase has only recently been recognized and this has prompted several studies to be carried out, mainly for the following reasons.\nFirst, the enzyme may be responsible for the hydrolysis of the dietary sphingomyelin occurring substantially in milk, eggs, meat and fish. Second, this enzyme may regulate cholesterol absorption. Third, the presence of alkaline SMase along the intestinal tract and its selective decrease detected in colorectal carcinoma suggests that this enzyme plays a role in intestinal carcinogenesis, since under physiological conditions, it stimulates apoptosis and protects the intestinal mucosa against carcinogenesis.\nPrevious studies have also shown that, under physiological conditions, alkaline SMase is dissociated by bile salts from intestinal mucosal membrane to the lumen. However, under pathological conditions, whereby bile salt concentration is abnormally increased, the dissociation of alkaline SMase by bile salts may exceed the biosynthesis of the enzyme, resulting in a low level of activity of alkaline SMase in the mucosa, and an abnormally increased excretion of the enzyme in the feces or in biological fluids, i.e., bile. Consequently, the excess of alkaline SMase excreted in the stools or in biological fluids over normal, basal values, may be interpreted as a valuable diagnostic marker for colon rectal carcinoma and familial adenomatous polyposis, hence; the need of a reliable assay for detecting alkaline SMase in the stools or in biological fluids of patients likely to be suffering from the aforesaid pathologies of the intestinal tract.\nIn addition, some bacteria strains (e.g., Streptococcus htermophilus or Lactobacilli) contain high levels of SMase, and the assessment of alkaline SMase may provide a method to evaluate changes in the number of said bacteria, i.e., after a treatment with probiotics or/and probiotic-based products.\nPrevious methods for assaying alkaline SMase are already known. The activity of the SMases can be determined either in vivo through cell labeled with a radioactive precursor of SM and then determining the labeling product levels or in vitro using radiolabeled SM or a chromogenic analog of SM or colored and fluorescent derivatives of neutral SM.\nThese known commonly used assays are not entirely satisfactory since they are potentially very hazardous insofar as they are radioactive assays and less sensitive than a fluorometric assay."}
-{"text": "All primary headaches are in need of better treatments. The most important primary headaches (i.e. independent disorders that are not caused by another disease) are migraine, tension-type headache and cluster headache. Migraine has a prevalence of 10% in the general population with a lifetime prevalence of 13% in men and 33% in women. Migraine is a highly disabling disease with high personal and social costs. To date, the precise mechanisms underlying the pathophysiology of migraine have remained elusive. Migraine strikes people during what are expected to be their most productive years: between ages 20 and 40 for most women, with a slightly higher age range for men. Migraine is typically characterized by unilateral onset of head pain, severe progressive intensity of pain, throbbing or pounding, and interference with the person's routine activities. Accompanying symptoms of photophobia (sensitivity to light) or phonosensitivity (intolerance to noise), as well as nausea and/or vomiting, are common, and often leads to the inability to perform daily tasks. A large portion of people with migraine often have no accompanying pain, their predominant symptom instead being vertigo (a spinning sensation) or dizziness/disequilibrium (balance loss), mental confusion, disorientation, dysarthria, visual distortion or altered visual clarity, or extremity paresis. Patients with migraine associated vertigo (MAV) are often seen by audiologists and vestibular rehabilitation therapists for evaluation and treatment. Because the exact mechanisms of migraine are still not completely understood, the management of migraine dizziness presently includes a combination of medications, vestibular rehabilitation, and lifestyle modifications that include limitation of risk factors associated with migraine (those related to diet, sleep, stress, exercise, and environmental factors).\nMigraine is a disease associated with increased synthesis and release of calcitonin gene related peptide (CGRP) and a migraine attack can be blocked with CGRP antagonists. The actual pain is generated by nociceptors of trigeminal nerve endings in the dura. Low serotonin levels may sensitize the nociceptors of trigeminal neurons. Triptans and ergotamins, which decrease serotonin, are associated with relief of acute pain. In contrast, tricyclic antidepressants and selective serotonin and noradrenaline reuptake inhibitors, which are associated with increases in serotonin, are utilized for migraine prevention.\nMigraine attacks can be triggered by intrinsic cerebral factors (e.g. calcitonin gene related peptide (CGRP) release), nitric oxide like tri-nitroglycerine, corticotrophin releasing hormone (stress), pro-inflammatory cytokines, and degranulation of mast cells located in the dura. While migraine has a genetic background, twin studies reveal that the cause of a majority of migraines appears to be due to environmental factors. The mechanism of triggering migraine is, however, still not understood. The cause of both migraines and chronic dizziness has eluded investigators for centuries and it therefore presents a truly long felt but unsolved mystery as to its causation and treatment."}
-{"text": "Clinical studies and practice have shown that providing a reduced pressure in proximity to a tissue site augments and accelerates the growth of new tissue at the tissue site. The applications of this phenomenon are numerous, but application of reduced pressure has been particularly successful in treating wounds. This treatment (frequently referred to in the medical community as \u201cnegative pressure wound therapy,\u201d \u201creduced pressure therapy,\u201d or \u201cvacuum therapy\u201d) provides a number of benefits, including faster healing, and increased formulation of granulation tissue."}
-{"text": "This description relates to operation of security and surveillance systems.\nIt is common for businesses and homeowners to have a security system for detecting alarm conditions at their premises and signaling the conditions to a monitoring station or to authorized users of the security system. Security systems often include an intrusion detection panel that is electrically or wirelessly connected to a variety of sensors. Those sensors types typically include motion detectors, cameras, and proximity sensors (used to determine whether a door or window has been opened)."}
-{"text": "1. Field of the Invention\nThe invention relates to a tool saver for use in association with oil or gas wells or the like. More specifically, the invention relates to a tool saver which prevents tools, when withdrawn from a borehole of a well, from falling back into the well in the event of a break in the wire line which lowers and raises the tools.\n2. Description of Prior Art\nAs is known in the art, wire line tools for wells, such as, for example, sand line swabbing equipment, sinker bars, etc., are lowered into and raised out of wells through the well bores of the wells by wire lines. There is a constant danger that, after the tool is raised out of the well, the wire line will break and the tool will fall uncontrollably back into the well. The tool will then have to be retrieved at great cost and inconvenience.\nVarious devices are known for gripping tools or pipes as they are lowered into wells. Examples of such devices are taught in, for example, U.S. Pat. Nos. 1,422,289, Moody, Jul. 11, 1922, 4,076,337, Childress, Feb. 28, 1978, 4,096,608, Lagerstedt, Jun. 27, 1978 and 4,898,238, Grantom, Feb. 6, 1990.\nThe '289 patent teaches a device for gripping casings or sections of pipes as they are being raised or lowered into a well. Specifically, the device will hold one section while another section is being coupled or uncoupled.\nThe drill steel holder of the Childress patent, U.S. Pat. No. 4,076,337, comprises levers 26 and 28 (see FIGS. 3 and 4) having outboard ends connected to a RAM 40. The RAM 40 can draw the outboard ends towards each other or force them apart from each other.\nThe inboard ends of the levers 26 and 28 comprise rollers 50. When the inboard ends of the levers are drawn towards each other, as shown in FIG. 3, the rollers will embrace a tool, for example a steel drill 68, to prevent the vertical movement of the steel drill.\nJaws 24 and 25 of U.S. Pat. No. 4,096,608, (having inserts 28 and 29) are forced towards each other by a spring 49 to capture a drill pipe 19 as shown in FIG. 2 of the drawings. As seen in FIG. 1, this will prevent the drill pipe 19 from moving in a vertical direction.\nIn U.S. Pat. No. 4,898,238, Grantom, a stationary jaw 200 acts together with a gripping, or movable, jaw 300 to capture a pipe or the like which is inserted into a well. As seen in FIG. 1, the movable jaw 300 is moved by the action of threaded shaft 320 in threaded block 310.\nAs is quite clear, none of the references address themselves to the problem as above described."}
-{"text": "This invention relates generally to apparatus and methods for making induced polarization and resistivity measurements in a borehole.\nIn the study of geological formations it is known that induced polarization and resistivity measurements can be made by means of surface-located equipment. See, for example, U.S. Pat. No. 4,295,056. Through the use of such surface equipment, induced polarization measurements which may include phase angle anomalies can be obtained. The phase angle anomalies can be created by a shallow subsurface polarizable material which is of significance to the oil and gas industry if its composition can be geochemically tested and found to contain hydrocarbons which have seeped from a deeper reservoir. However, before the material can be tested, its location must be pinpointed. Therefore, there is the need for an apparatus and a method which utilize electrodes in a borehole to make downhole induced polarization and resistivity measurements in the vicinity of the formations which are creating the anomalies so their locations can be accurately determined. Once located, the formations can be geochemically tested to determine if hydrocarbons are present (see, for example, U.S. Pat. No. 4,360,359). From such information one can tell whether hydrocarbons are seeping from a reservoir, for example. The need for an apparatus and a method for making such downhole measurements has been recognized. See, for example, Complex Formation Resistivity--The Forgotten Half of the Resistivity Log by D. D. Snyder, R. H. Merkel, and J. T. Williams and Induced Polarization Measurements In and Around Boreholes by D. D. Snyder and R. H. Merkel.\nIt has been observed that some of these anomalies arise from shallow or thin formations, thereby making it desirable to utilize an apparatus which can make measurements of sufficient resolution for such thin formations. This desirability indicates that a short array spacing of the electrodes be used, which short array spacing corresponds to the size of the thin layer; however, for such a short array spacing of the type needed to make the induced polarization and resistivity measurements, the measurements made thereby would be adversely affected by the electrical properties of the drilling mud or other fluid surrounding the electrodes at their measurement locations in the borehole. That is, a short array spacing would cause the apparatus to read the electrical properties of the drilling mud or other fluid rather than the electrical properties of the desired formation.\nThe borehole fluid can have a substantial effect on the phase readings made by such an apparatus because the borehole fluid may be a drilling mud typically containing highly polarizable montmorillonite clay, for example. Additionally, it is difficult, if not virtually impossible, to determine ahead of time exactly what the effect of the borehole fluid will be. Still further, a drilling mud tends to be electrochemically unstable so that it is difficult to determine at any given time what the electrical properties of the mud are. Furthermore, the depth of invasion of the mud or other borehole fluid into the rock formation is highly variable. Therefore, there is the need for the apparatus and method to be constructed and performed so that several suitable measurements with variable length arrays can be made and compared to derive a reasonable estimate of the electrical properties of thin formations. This reasonable estimate can be further refined by suitable modeling of the formation surrounding the borehole.\nIt is also desirable that the apparatus and method which meet the foregoing needs be usable with commercially available interconnection equipment presently found in the oil and gas industry. This commercially available equipment includes electrical logging cable comprising parallel electrical conductors. Such commercially available cable is susceptible to significant electromagnetic coupling between the various parallel conductors. This electromagnetic coupling adversely affects the induced polarization measurements.\nOne way to avoid the electromagnetic coupling is to use separately twisted pairs of wires; however, such twisted pairs of wires are not standard equipment in the oil and gas industry and thus their use would require special orders with possibly long lead times and high costs. Additionally, separately twisted wires would make for a thick cable which would reduce the amount of cable which could be retained on available winches.\nSuch a problem of electromagnetic coupling could also be overcome by placing certain equipment in the implement which is to be used in the borehole and telemetering the readings to the surface. A shortcoming of placing additional equipment in the borehole is that it requires packaging the equipment for use in the extreme pressures and temperatures found in the borehole. Furthermore, having the equipment downhole would limit the ability to control it, which ability would be present if the equipment were to be located at the surface. Telemetering signals to the surface would limit the dynamic range and accuracy of the measurements made in the borehole. Therefore, there is the need for the apparatus and method to be constructed and performed so that such adverse electromagnetic coupling is eliminated without preventing the apparatus and method from being used with or implemented by commercially available equipment.\nIt is known that such induced polarization and resistivity measurements as are referred to hereinabove are made by introducing a current into the formation and then detecting a voltage produced in the formation by the current. It has been discovered that it is important to control the density of the current because if the current density is too large, the induced polarization response will be non-linear with respect to the current and thus will not be representative of the response that a surface induced polarization survey would obtain. Therefore, the apparatus and method which meet the aforementioned needs should also meet the needs of maintaining an accurately controlled current output and utilizing a low-noise receiver which can accommodate very small transmitted currents.\nTo insure that accurate measurements are made, the apparatus should be constructed so that electrical leakage in the portion of the apparatus which is to be used in the borehole is reduced. It has been found that the phase measurement which is indicative of the induced polarization response is much more sensitive to minor amounts of electrical leakage than are conventional magnitude-only resistivity measurements. To prevent such electrical leakage, there is the need for the apparatus to be constructed so that there is substantially a zero pressure differential between the interior of the downhole implement and the borehole pressure to thereby preclude the flow of borehole fluid into the downhole implement.\nTo further insure the accuracy, the desired method is to be performed discretely whereby discrete, rather than continuous, measurements of the induced polarization and resistivity are made. Discrete measurements preclude the introduction of noise and instability which can occur by means of the movement of the downhole implement through the borehole fluid as measurements are taken. Additionally, the preferred current used to obtain such induced polarization and resistivity measurements is at a very low frequency so that if a continuous logging run were made, it could cover a substantial distance before one cycle of the current were transmitted.\nMore generally, there is the need for an apparatus and a method for obtaining measurements which can be used not only to qualitatively look for anomalies, but also to quantitatively understand surface-obtained induced polarization data, other surface geological and geophysical data and downhole electrical properties."}
-{"text": "The systemic administration of therapeutic agents treats the body as a whole even though the disease to be treated may be localized. In some cases of localized disease, systemic administration may not be desirable because the drug agents may have unwanted effects on parts of the body which are not to be treated or because treatment of the diseased part of the body requires a high concentration of a drug agent that may not be achievable by systemic administration.\nIt is therefore often desirable to administer therapeutic agents to only localized sites within the body. Common examples of where this is needed include cases of localized disease (e.g., coronary heart disease) and occlusions, lesions, or other disease in body lumens. Several devices and methods for localized drug delivery are known. In one example, such devices are drug delivery balloons, and methods of their use include the steps of coating a balloon attached to a balloon catheter with a drug and a carrier matrix, inserting the catheter into a blood vessel, tracking the balloon to a desired location, and expanding the balloon against the surrounding tissue to transfer the drug locally at the intended treatment site.\nOne of the potential drawbacks to localized drug delivery is the possibility of premature or unintended release of the drug, the carrier matrix, and/or the drug/carrier matrix combination. This may occur during tracking and placement at the treatment site of a drug delivery device and post delivery as the device is withdrawn from the body. Such unintended release may result from drug diffusion, device contact with areas proximate the treatment site, or washing of the drug from the surface of the delivery device due to blood flow. This is of particular concern when the device comprises a therapeutic agent of a type or dosage not intended to be released to tissue or blood outside the treatment site.\nDrugs or coating components shed in this unwanted fashion may be in particulate form or may be in solution. The downstream release of undesirable particles is known as \u201cparticulation\u201d. For example, particulation of large particles can create problems such as ischemia in tissues, especially in tissues supplied by small diameter vessels. Furthermore, the resulting effects of biodistribution of such particles are not well understood and may result in adverse effects.\nIn view of the potential drawbacks to current, localized drug delivery, there exists a need for devices and methods that allow for controlled, localized delivery of drug agents to specific treatment sites within a mammalian body that avoids premature or unintended drug release away from the intended treatment site."}
-{"text": "The present invention relates generally to wearable computers and, more specifically, to methods and arrangements for facilitating the connection of auxiliary components therewith.\nAs miniaturization of computer products continues, the problem of adequate space for connections of input and output devices (I/O) continues to grow. The minimum size of many computer products are now being defined by the number and type of connectors used, especially on the highly integrated devices of the type used in wearable computer technology. In other words, since it is the case that a minimum amount of space must typically be allocated for accommodating connectors or connecting interfaces in the environment of a wearable computer, it has been found that further efforts towards miniaturization in the context of wearable computers has been hampered.\nA need has thus been recognized in connection with facilitating further miniaturization in the context of wearable computers by reducing the amount of space associated with connectors, connection interfaces or other connecting elements.\nIn accordance with at least one presently preferred embodiment of the present invention, a single docking-type connector is provided on a computer system unit, and distributes the I/O connectors required for a given task, into the harness (wearing device) of the computer. The docking connector preferably has enough pins to carry all of the I/O signals, power and grounds for the computer.\nSuccessful wearable computers will offer flexible configurations of I/O and wearing possibilities. Because the system unit has all of the possible I/O signals, different types of harnesses (as variously illustrated by way of example in the drawings) allow for many configurations. The harnesses with the cabling are relatively inexpensive compared to the system unit.\nGenerally, in accordance with at least one presently preferred embodiment of the present invention, problems relating to the space taken up by connectors are solved by essentially exchanging the multiple connectors required for I/O on a system unit for a single docking connector. By redistributing the individual connectors into various harnesses and docking stations, the system unit is greatly reduced in size and can achieve. attendant advantages as discussed herein.\nIn one aspect, the present invention provides a connector system for a computer, the connector system comprising: a docking connector which interfaces with a system unit; a plurality of leads in communication with the docking connector; and a plurality of connecting elements associated with the plurality of leads, the connecting elements each being adapted to interface with at least one external component.\nIn another aspect, the present invention provides a harness system for supporting a portable computer, the harness system comprising: a harness which accommodates a system unit; a docking connector which interfaces with the system unit; a plurality of leads in communication with the docking connector; and a plurality of connecting elements associated with the plurality of leads, the connecting elements each being adapted to interface with at least one external component.\nFor a better understanding of the present invention, together with other and further features and advantages thereof, reference is made to the following description, taken in conjunction with the accompanying drawings, and the scope of the invention will be pointed out in the appended claims."}
-{"text": "1. Field of the Invention\nThe invention relates to devices and methods for treating incontinence. More specifically, particular embodiments of the invention use a variable control vacuum and/or vibratory device or other device to apply suction or other force to the male or female genital region, e.g. the region of the female urethral opening, the clitoral region and/or the vaginal region, or the male penile region.\n2. Description of Related Art\nA. Incontinence\nUrinary incontinence is a significant clinical problem and a major source of disability and dependency. Although urinary incontinence is particularly prevalent in the elderly, it affects all age groups. Millions of male and female adult Americans suffer from urinary incontinence, including at least one-half of all nursing home residents. The monetary costs of managing urinary incontinence, and the psychological costs associated with the problem, are great. Fecal incontinence presents similar concerns.\nThe most frequently occurring types of urinary incontinence are stress incontinence, urge incontinence, overflow incontinence, and mixed incontinence. Stress incontinence is a common form of incontinence in women. Intra-abdominal pressure exceeds urethral pressure upon coughing, sneezing, laughing, lifting, or like activity, causing leakage of urine. Physical changes associated with pregnancy, childbirth, and menopause, for example, are known to cause stress incontinence.\nUrge incontinence occurs when a patient loses urine while suddenly feeling the urge to urinate. The patient is unable to inhibit the flow of urine long enough to reach the toilet. Inappropriate bladder contractions are the most common cause of urge incontinence, and may occur in connection with central nervous system lesions, urinary infection, or bladder tumors, to name several examples.\nOverflow incontinence occurs when the bladder is unable to empty normally. Weak bladder muscles, caused e.g. by nerve damage from diabetes, or a blocked urethra, caused e.g. by tumors or urinary stones, are among the more common causes of overflow incontinence. Frequency or urgency involves the need or urge to urinate on an excessively frequent or habitual basis. Combinations of these and other types of incontinence, e.g. stress incontinence and urge incontinence, are often called mixed incontinence.\nMany options are available to treat incontinence in its various forms, including Kegel exercises, electrical stimulation, biofeedback, timed voiding or bladder training, medications, pessaries, implants, invasive or minimally invasive surgery, catheterization, and other methods and devices. Kegel exercises, or pelvic floor muscle exercises, for example, are intended to strengthen the muscles supporting the urethra, bladder, uterus and rectum. Weak pelvic floor muscles can contribute to urinary incontinence problems; hence strengthening those muscles can tend to alleviate those problems.\nAccording to several known surgical treatment methods, one or more stimulation systems, including electrodes controlled by an electronic control system, are implanted to electrically stimulate nerves controlling external sphincter and bladder functions. The electronic control system directs the administration of electric pulses to stimulate the nerves appropriately. Note U.S. Pat. No. 4,771,779, for example, which is incorporated herein by reference. U.S. Pat. No. 6,055,456, also incorporated herein by reference, shows an implantable medical lead for stimulation of the sacral nerves to treat incontinence.\nAlthough treatments requiring surgical intervention may be the preferred and most effective treatment mode in some situations, surgical intervention may be too extreme a measure in other situations. In some cases, surgical procedures to treat incontinence actually have a relatively low success rate; in many cases such procedures are irreversible. Additionally, a patient may hesitate to proceed with a surgical option, and/or a patient's physical condition may make surgical intervention inappropriate. Surgery may be inappropriate for pregnant patients, for example, or those of advanced age. Similarly, pharmacological treatment options may cause undesirable side effects and/or interactions with other medications. Non-surgical treatments, for example exercises or bladder training, may demand too high a degree of patient compliance or effort and thus may be resisted or otherwise ineffective.\nA need exists, therefore, to treat urinary and/or fecal incontinence in a non-invasive, non-pharmacological manner that is less likely to meet with patient resistance or cause physical trauma or side effects.\nB. Urogenital Anatomy\nThe female urethra passes just anterior to the vagina, and the external female urethral orifice is just posterior to the clitoris. The external urethral sphincter and the compressor urethrae compress the urethra and serve as functional sphincters. The pudendal nerve, which is derived from the anterior divisions of the ventral rami of S2 through S4, is the chief nerve of the perineum, providing sensation to the genitalia including the clitoris and controlling the motor function of the external urethral sphincter and the external anal sphincter. In the male, the pudendal nerve provides sensation from the penis and scrotum and provides motor control to the pelvic floor, including the anal sphincter.\nC. Clitoral Anatomy\nThe clitoris in the human female consists of a cylindrical, erectile organ composed of three parts: the outermost glans or head, the middle corpus or body, and the innermost crura. The glans of the clitoris is visualized as it emerges from the labia minora, which bifurcates to form the upper prepuce anteriorly and the lower frenulum posteriorly. The body of the clitoris consists of two paired corpora cavernosa of about 2.5 cm in length. The body extends under the skin at the corona to the crura. The two crura of the clitoris, formed from the separation of the most proximal portions of the corpora in the perineum, attach bilaterally to the undersurface of the symphysis pubis at the ischiopubic rami.\nA fibrous tunica albuginea ensheathes each corporal body made up of lacunar space sinusoids surrounded by trabecula of the vascular smooth muscle and collagen connective tissue. No retractor clitoridis muscle exists in humans as it does in other animals such as cattle and sheep, however a supporting suspensory ligament does hold the clitoris in the introital region.\nThe main arterial supply to the clitoris is from the ilio-hypogastric-pudendal arterial bed. The internal pudendal artery is the last anterior branch off the internal iliac artery. Distally, the internal pudendal artery traverses Alcock's canal, then terminates as it supplies the inferior rectal and perineal artery which supply the labia. The common clitoral artery continues to the clitoris. This artery bifurcates into a dorsal clitoral artery and a cavernosal clitoral artery.\nIn the normal female, autonomic efferent innervation of the clitoris passes from the pelvic and hypogastric nerves to the clitoris. Pelvic nerve stimulation results in clitoral smooth muscle relaxation and arterial smooth muscle dilation, causing an increase in clitoral cavernosal artery inflow and an increase in clitoral intracavernous pressure, which lead to tumescence and extrusion of the glans clitoris.\nThe clitoris has a dense collection of Pacinian corpuscles, Meissner's corpuscles, Merckel tactile disks, and free nerve endings. These sensory afferent nerves pass through the dorsal clitoral nerve to the pudendal nerve and into the sacral nerve roots (S2, S3, S4). (The male anatomy is completely homologous.) The presence of at least one somatic reflex arc is demonstrated by the bulbo-cavernosal reflex; squeezing the clitoris causes the anal sphincter to contract\u2014the so-called \u201canal wink.\u201d\nD. Female Sexual Dysfunction\nClitoral erectile insufficiency or reduced clitoral arterial flow may be caused by atherosclerosis, diabetes, or age-related causes, among other factors. Reduced clitoral arterial flow may lead to fibrosis of the clitoral cavernosa and reduced clitoral physiological function. In an animal model, Park et al. demonstrated that significant collagen synthesis occurs when the arterial inflow to the clitoris is compromised. This work demonstrated the importance of maintaining arterial flow to the clitoris to prevent collagen synthesis and fibrosis on the smooth muscle. See Park, K., et al., Vasculogenic Female Sexual Dysfunction: The Hemodynamic Basis for Vaginal Engorgement Insufficiency and Clitoral Erectile Insufficiency, IJIR, 9:27-37, 1997.\nIt is believed that the difficulty or inability to achieve clitoral tumescence may be related to and associated with other symptoms of female sexual arousal disorder. According to the International Consensus Report on Female Sexual Dysfunction, Female Sexual Arousal Disorder (FSAD) is defined as the persistent or recurrent inability to attain or maintain adequate genital lubrication or swelling responses resulting in personal distress. FSAD may be expressed as a lack of subjective excitement or lack of genital (lubrication/swelling) or other somatic responses (AFUD Consensus Report of FSD, 1998).\nA non-pharmacological approach to treatment that causes blood flow and engorgement, thereby applying a stimulus to the sensory nerve endings in the clitoris, periurethral area, and/or genital area, would be very beneficial to a large group of women complaining of FSAD."}
-{"text": "The present invention relates to a focusing unit for controlling the focusing voltage of a cathode ray tube within a TV receiver by appropriately controlling the high voltage of the FBT (flyback transformer) and, more particularly, to a dynamic focusing unit capable of controlling the peripheral focus in the up and down and the right and left dimensions of the TV screen correctly and over a wider range.\nGenerally, the focusing unit is a device for focusing the electron beam, being an electrostatic lens formed between anode electrode and cathode electrodes by applying pertinent voltage to the focusing electrodes such that a clear image is displayed on a fluorescent screen by the electron beam emitted from the electron gun of the picture tube.\nFIG. 8 is a schematic structural diagram of a picture tube with a spherical screen controlled by the focusing circuit of Korean Utility Model official publication gazette No. 89-8220, and wherein the electrons emitted from cathode 100 of the electron gun are made into an electron beam being radiated to a beam focusing point on a fluorescent screen 130. For focusing this electron beam, focusing electrode 110 is provided within the interior of the picture tube and then pertinent voltage is applied from the exterior. Thereby, the electron beam is focused by the electrostatic field of focusing electrode 110.\nHowever, although the electron beam focused upon the fluorescent screen 130 of the aforementioned spherical picture tube does not produce deviation, in the case when the fluorescent screen 130 is flat as the picture screen 140, a focusing deviation 150 is produced at the peripheral regions of the flat screen.\nAccordingly, since the conventional focusing device was a device designed for a spherical cathode ray tube, a focusing deviation has been produced at the peripheral region of the flat cathode ray tube.\nIn order to solve the above-described defect, a device has been developed in which focusing is correctly controlled at peripheral regions of cathode ray tubes by utilizing the horizontal flyback pulse signal and vertical output signal.\nAs shown in FIG. 9, a focusing device for controlling the electron beam to be focused on the flat screen of cathode ray tube 300 in response to the signal applied to focusing electrode 110 comprises a vertical signal integrating circuit 10 for integrating vertical pulse signal S.sub.1, thereby converting to a vertical parabolic waveform signal S.sub.2, a vertical amplifying circuit 20 for variably amplifying the output signal of said vertical signal integrating circuit 10, thereby controlling the up and down focusing on the screen, a horizontal signal integrating circuit 30 for integrating the horizontal pulse signal S.sub.4 inputted after being extracted from the flyback transformer, thereby converting to a horizontal parabolic waveform signal S.sub.5, a buffer circuit 40 for variably current amplifying the output signal of said horizontal signal integrating circuit, thereby controlling the right and left focusing on the screen, and an inversion amplifying circuit 50 for inversely amplifying the mixed vertical and horizontal parabolic waveform signals S.sub.6. The signal S.sub.7 outputted from said inversion amplifying circuit 50 is applied to focusing electrode 110 of picture tube 300 and also to anode 160 of picture tube 300. These elements are arranged so that the phenomenon of deviating focus at peripheral regions of the screen is prevented. The focusing unit serves a function which controls appropriately the high voltage coming out of the flyback transformer (FBT) thereby controlling the focus of the cathode ray tube within the TV receiving set.\nHowever, there are many kinds of cathode ray tubes and different focusing voltages are required by each kind. In order to feed focusing voltage of the cathode ray tube for each kind, the wider the variable range of the voltage of the focusing unit is, the better.\nA conventional focusing unit is shown in FIGS. 6 and 7 and comprises a knob 208 passing through hole 201 in the rectangular main body 200 of the focusing unit having annular ring 204 protruded with two rotation limiting protrusions 202, 203 extending inward, contact terminal fixing means 206 having a horseshoe shape to receive variable contact pin 205, knob body 210 structured with variable shaft means 207 formed integrally with said fixing means 206 as well as knob 208, and ceramic base board 230 provided with input and output voltage terminals 231, 232, circular printed portion 233 and bent printed portion 234.\nHowever, in this conventional focusing unit, since variable contact pin 205 is moved along circular printed portion 233 around knob 208, the extent of the output voltage terminal 231 from the input voltage 25 terminal 232 is limited within a narrow range.\nTherefore, since it has a simple structure that the first focusing voltage variable range 240 cannot cover by less than 10%-15% and also the second focusing voltage variable range 250 cannot cover by less than 2%-3%, there has been a problem of achieving a satisfactory variable voltage focusing range."}
-{"text": "The present invention relates to liquid crystal display devices, and, more particularly, to liquid crystal display devices of the lateral electric field type.\nA liquid crystal display device of the lateral electric field-driven type comprises a pair of transparent substrates disposed so as to be opposed to each other and to be spaced by a liquid crystal (LC) layer sealed therebetween, and a pixel electrode and a counter electrode are spaced apart from the pixel electrode formed in each of the pixel regions on a surface at the liquid crystal layer side of one of the pair of transparent substrates. A pixel electrode and a counter electrode disposed at each of the pixel regions generates an electric field therebetween so as to control the optical transmissivity of the liquid crystals in the vicinity of the pixel region. Since the liquid crystal display device of the lateral electric field type generates an electric field having a major component extending substantially parallel to a main surface of the transparent substrate between the pixel electrode and the counter electrode for switching the orientations of the liquid crystal molecules, the lateral electric field scheme is also referred to as the xe2x80x9cIn-Plane-Switching (IPS) schemexe2x80x9d also.\nAn active-matrix liquid crystal display device employing the lateral electric field scheme has a plurality of scan signal lines and a plurality of counter voltage signal lines both of which extend in an xe2x80x9cxxe2x80x9d direction (row direction), while being juxtaposed in parallel in a xe2x80x9cyxe2x80x9d (column) direction (transverse to the x direction), and a plurality of video signal lines which extend in the y direction, while being juxtaposed in the x direction on the surface at the liquid crystal layer side of the one of the pair of transparent substrates. When the scan signal lines, the counter voltage signal lines and the video signal lines are disposed on a surface of one of the pair of transparent substrates, each of the pixel regions is defined as an area surrounded by one of the scan signal lines, one of the counter voltage signal lines disposed adjacent to the scan signal line, and a pair of the image signal lines disposed adjacent to one another. Each of the pixel regions has a switching element, a pixel electrode, and a counter electrode. The switching element operates so as to be turned on by a scanning signal supplied by one of the scanning signal lines. The switching element, when turned on, supplies a video signal being transmitted by one of the video signal lines thereto to the pixel electrode in a pixel region corresponding thereto. A counter voltage for generating a potential difference between the counter electrode and the pixel electrode in the pixel region is transmitted through one of the counter voltage signal lines and is applied to the counter electrode.\nA liquid crystal display device having the structure described above is characterized by the fact that it produces an image having a higher contrast ratio and exhibits a remarkably wider viewing angle.\nThe lateral electric field-type liquid crystal display device referred to above is described in detail in such documents as Japanese Patent Application Laid Open Hei 05-505247 (No. 505247/1993), Japanese Patent Publication Sho 63-021907 (No. 021907/1988), and Japanese Patent Application Laid Open Hei 06-160878 (No. 160878/1994).\nThe lateral electric field-type liquid crystal display device having the structural features described above utilizes electrodes (e.g. an pixel electrode and a counter electrode) formed of thin conductive films of several thousands Angstrom (xc3x85) (thin film electrodes, hereinafter) and generates an electric field substantially in parallel with a surface of one of the pair of transparent substrates. However, it is difficult to reproduce the aforementioned electric field sufficiently and accurately in the liquid crystal layer lying between the thin film electrodes.\nIn addition to the aforementioned difficulty, the aforementioned electric field being generated between the pixel electrode and the counter electrode in a pixel corresponding thereto is affected by another electric field appearing from the video signal line adjacent to the pixel.\nOn the other hand, a black matrix layer (a light shielding film) is formed on a surface of another of the pair of transparent substrates confronting the one of the pair of transparent substrates on which the pixel electrode and the counter electrode are formed. For instance, the black matrix layer is formed on the surface of the substrate which confronts the liquid crystal layer. The black matrix layer may be formed of an insulating material or the like having such a high resistivity that the electric field between the pixel electrode and the counter electrode is generated to be as strong as possible.\nThe black matrix, being formed of an insulating material or the like, tends to have a higher optical transmissivity than that of metals so that a problem of deterioration of the contrast ratio occurs in an image displaying operation of the liquid crystal display device.\nThe present invention has been made in view of the technical background mentioned above, and an object of the present invention is to provide an improved liquid crystal display device which is suitable for generating an electric field sufficiently between the pixel electrode and the counter electrode as to control the orientations of the liquid crystal molecules lying therebetween in spite of utilizing a simple structure.\nAnother object of the present invention is to improve the contrast ratio of an image displayed by the liquid crystal display device.\nSome representative aspects and features of the present invention as disclosed herein will be briefly summarized as follows.\nThe present invention provides a liquid crystal display device comprising: a pair of substrates being superposed on one another with a predetermined space therebetween in which a liquid crystal layer is sealed; at least one video signal line being formed on a surface of one of the pair of substrates confronting the liquid crystal layer; and at least one pixel being formed on the surface of the one of the pair of substrates and having at least one pixel electrode to which a signal is supplied through the at least one video signal line and at least two counter electrodes being spaced from the at least one pixel electrode, one of the at least two counter electrodes being disposed adjacent to the at least one video signal line. The display device is characterized in that: (1) the at least one video signal line and the one of the at least two counter electrodes are spaced from one another by at least two insulating regions disposed along the surface, and (2) one of the at least two insulating regions comprises a material having a higher permittivity than that of the other of the at least two insulating regions.\nThe present invention also provides another liquid crystal display device comprising: a pair of substrates being disposed opposite to one another a liquid crystal layer being sealed therebetween; at least one video signal line formed on a surface of one of the pair of substrates confronting the liquid crystal layer; and at least one pixel being formed on the surface of the one of the pair of substrates, being arranged adjacent to the at least one video signal line, and having at least one pixel electrode and at least one counter electrode, so as to apply an electric field to the liquid crystal layer therebetween. The display device is characterized in that: (1) at least a surface of the at least one video signal line is covered with an insulating material having a permittivity (a dielectric constant) higher than 6.7.\nThese and other objects, features and advantages of the present invention will become more apparent from the following description when taken in conjunction with the accompanying drawings."}
-{"text": "Of all of the applications where computer vision is used, face detection presents an extremely difficult challenge. For example, in images acquired by surveillance cameras, the lighting of a scene is usually poor and uncontrollable, and the cameras are of low quality and usually distant from potentially important parts of the scene. Significant events are unpredictable. Often, a significant event is people entering a scene. People are typically identified by their faces. The location and orientation of the faces in the scene are usually not controlled. In other words, the images to be analyzed are substantially unconstrained.\nIn applications where a large number of images are to be analyzed, or the analysis needs to be performed in real-time, the number of images that actually include faces is usually very small. Moreover, in most images that do include a face, the amount of the image that is a face is also very small. Therefore, it is desirable to rapidly eliminate images or portions of images that do not include a face, thus more processing can be devoted to portions of images that potentially contain a face.\nFace detection has a long and rich history. Some techniques use neural network systems, see Rowley et al., \u201cNeural network-based face detection,\u201d IEEE Patt. Anal. Mach. Intell., Vol. 20, pages 22\u201338, 1998. Others use Bayesian statistical models, see Schneiderman et al., \u201cA statistical method for 3D object detection applied to faces and cars,\u201d Computer Vision and Pattern Recognition, 2000. While neural network systems are fast and work well, Bayesian systems have better detection rates at the expense of longer processing time.\nU.S. Pat. No. 5,710,833, \u201cDetection, recognition and coding of complex objects using probabilistic eigenspace analysis\u201d issued to Moghaddam, et al. on Jan. 20, 1998, describes a system for detecting instances of a selected object or object feature, e.g., faces, in a digitally represented scene. Their method utilizes analysis of probability densities to determine whether an input image, or portion thereof, represents such an instance.\nU.S. Pat. No. 6,337,927, \u201cApproximated invariant method for pattern detection\u201d issued to Elad et al. on Jan. 8, 2002, describes a system and method for classifying input feature vectors into one of two classes, a target class (faces) and a non-target class (non-faces). The system utilizes iterative rejection stages to first label the input vectors that belong in the non-target class in order to identify the remaining non-labeled input vectors as belonging in the target class. The operation of their system is divided into an off-line (training) procedure and an online (actual classification) procedure.\nDuring the off-line procedure, projection vectors and their corresponding threshold values that will be used during the on-line procedure are computed using a training set of sample non-targets and sample targets. Each projection vector facilitates identification of a large portion of the sample non-targets as belonging in the non-target class for a given set of sample targets and sample non-targets.\nDuring the on-line procedure, an input vector is successively projected onto each computed projection vector and compared with a pair of corresponding threshold values to determine whether the input vector is a non-target. If the input vector is not determined to be a non-target during the successive projection and thresholding, the input vector is classified as a target.\nAs a disadvantage, the features used by their method are limited to linear projections and upper and lower thresholds. Thus, outcomes of early features do not effect the outcomes of later tests. In other words, their tests can be reordered with no effect on the final result. In addition, because the classifier is a conjunction of linear projections followed by a threshold, the resulting classification function is limited to a convex region in the input space, which is a strong limitation that makes their system a weak classifier.\nA new framework for feature detection is described by Viola et al., in \u201cRapid Object Detection using a Boosted Cascade of Simple Features,\u201d Proceedings IEEE Conf. on Computer Vision and Pattern Recognition, 2001. They present three new insights: a set of image features which are both extremely efficient and effective for face analysis, a feature selection process based on Adaboost, and a cascaded architecture for learning and detection. Adaboost provides an effective learning algorithm and strong bounds on generalized performance, see Freund et al., \u201cA decision-theoretic generalization of on-line learning and an application to boosting,\u201d Computational Learning Theory, Eurocolt '95, pages 23\u201337. Springer-Verlag, 1995, Schapire et al., \u201cBoosting the margin: A new explanation for the effectiveness of voting methods,\u201d Proceedings of the Fourteenth International Conference on Machine Learning, 1997, Tieu et al., \u201cBoosting image retrieval,\u201d International Conference on Computer Vision, 2000.\nHowever, there are several problems with their cascaded approach. The simplest is that the cascade allows only for rejection thresholds, and does not have acceptance thresholds. This degrades performance when the distribution of acceptable and rejected examples are substantially equal.\nAnother disadvantage of their method is that later stages of the cascade ignore detailed information that was available during earlier stages of the cascade. If the early cascade stages accept features only marginally close to the threshold, then this useful information plays no role in later decisions to reject or accept the feature. For example, it is possible for multiple stages to accept a given patch because it is significantly above the rejection threshold, and then subsequently reject that patch as soon as it falls below the rejection threshold in later stages. In this case it is often better to accept the given patch.\nAnother disadvantage is their mechanism for setting the rejection thresholds. Their rejection thresholds are set for each cascade stage starting from the first stage. For any given task, their goal is to detect most of the objects, though it is impossible to detect all objects of a particular class. In the case of face detection, perhaps 90% to 95% of the faces are detected. The rejection threshold is set so that a very high percentage of the examples are detected. Therefore, in the early stages, one must be very conservative and not discard a potentially correct patch. For early stages in the cascade, 99% or more of the objects must be detected in order to ensure that valuable patches are not discarded.\nTheir process does not distinguish example patches that are easy to classify, and are eventually classified correctly, and those example patches, which are very difficult to classify and are eventually misclassified. If one could determine which example patches would eventually be misclassified, after the entire cascade, then those example patches could be eliminated early with no loss to the accuracy of the overall process. Because the rejection threshold setting process starts with the early stages of the cascade, these \u201cdifficult\u201d examples cannot be identified so the rejection threshold is set for all positive examples. This results in rejection thresholds which are lower and less effective in reducing processing time.\nTherefore, there is a need for an object detection system that improves upon the prior art."}
-{"text": "The subject matter hereof is related to the subject matter of U.S. Pat. No. 4,061,389 and is directed specifically to the improvement that promotes effective connection of a pull tool to a down hole object or fish and simple and effective disconnection of the pull tool from the article if such is desired.\nDuring drilling operations, well servicing operations and the like, objects may become lost or stuck within a well bore and these objects are typically referred to in the industry as \"fish.\" For example, a section of drill pipe or production tubing might become disconnected from a pipe string and it is then necessary to introduce a pull tool into the well bore, establish connection with the fish and then remove the fish. Many different types of pull tools, grap pling devices, spears, etc. have been developed for the purpose of conducting fishing operations in well bores. It is typical for wire line operations to be utilized for the purpose of introducing a pull tool into the well bore for fishing operations.\nIn many cases, down hole devices are positioned within casing and tubing strings for the purpose of controlling well operations. In most cases, such down hole devices are equipped with API Standard fishing necks which normally include an undercut shoulder to enable grappling by standard installation and retrieval tools. Quite often the API Standard fishing neck will become corroded or eroded by well conditions to the point that only a stub pipe is exposed without the usual undercut shoulder that is provided on the fishing neck. Connection between a wire line controlled fishing tool and a worn fishing neck may be accomplished by means of an overshot type grappling device generally defined by a collet structure having internal teeth that establish a gripping relation with the worn fishing neck. Although overshot type retrieving tools are successfully utilized in many cases, a common problem with such tools is the inability of the operator to achieve disconnection from the worn fishing neck in the event the tool is unable to accomplish an effective pulling operation. Occasionally, disconnection can be achieved by a substantial jarring or other violent mechanical movement, but, in some cases, disconnection of the pulling tool from the fish is extremely difficult if not impossible. In cases where the pulling tool or wire line is inadequate for the pulling operations that are required, it is necessary that the pulling tool be disconnected from the fish and replaced with a wire line controlled pulling tool of substantially greater pulling capacity. If a light-weight wire line pulling tool is unable to accomplish the pulling operation and becomes firmly fixed to the fish or other object to be pulled, retrieval of the fishing equipment itself obviously compounds the problem and adds materially to the expense of the service operation.\nAs mentioned above, application of shocks to the pulling tool will sometimes result in disconnection, thereby enabling service personnel to remove the wire line pulling tool and substitute a tool and wire line of substantially heavier gauge. Application of mechanical shocks to the down hole equipment of the well and the wire line tool itself can easily cause failure or excessive wear of one or more of these mechanical components. It is desirable, therefore, to provide down hole service equipment, such as wire line controlled pulling tools, that will effectively become interconnected with down hole objects such as wire line tools, fish, etc., and, in the event disconnection is necessary, will readily become disconnected upon simple mechanical movement of a wire line controlled pulling tool."}
-{"text": "The invention relates to marine propulsion devices. More particularly, the invention relates to outboard motors, and still more particularly to outboard motor tilt tubes.\nA conventional outboard motor, such as the one disclosed in U.S. Pat. No. 4,710,141, comprises a transom or stern bracket mounted on the transom of a boat, and a tilt tube which is supported by the transom bracket and which extends along a generally horizontal tilt axis. Nuts threaded onto the opposite ends of the tilt tube prevent axial movement of the tilt tube relative to the transom bracket. A swivel bracket is mounted on the tilt tube for pivotal movement relative thereto about the tilt axis, and a propulsion unit is mounted on the swivel bracket for pivotal movement relative thereto about a generally vertical steering axis. A steering ram slidably housed within the tilt tube has one end connected to the steering arm of the propulsion unit via a power steering system, and has the opposite end connected to the core of a steering cable assembly. The sheath of the steering cable assembly is fixed to one end of the tilt tube by a nut threaded onto the end of the tilt tube.\nA conventional tilt tube extends a certain distance beyond one side of the transom bracket, which distance is only long enough to enable a retaining nut to be threaded onto the end of the tilt tube, but not long enough to enable the sheath of a remote control cable to be connected to that end of the tilt tube. The other end of the tilt tube extends a sufficient distance beyond the transom bracket to enable a retaining nut to be threaded onto that end of the tilt tube and to enable a remote control cable sheath to be secured to that end of the tilt tube by an additional nut.\nAttention is directed to the following U.S. patents:\n______________________________________ Nolan 3,207,117 September 21, 1965 Baba 4,735,165 April 5, 1988 Hickham 4,815,994 March 28, 1989 ______________________________________"}
-{"text": "In some instances it is beneficial to utilize an apparatus comprising a compensating interferometer that modulates the phase of an input light signal to produce an output signal that is a superposition of an input signal and a phase shifted version of the input signal. One such instance is when a compensating interferometer is used to manipulate the output of a sensor array as described in U.S. Pat. No. 5,987,197. However, known compensating interferometers lack functionality that can be advantageous in processing optical signals."}
-{"text": "Rapid advances in electronics and communication technologies, driven by immense customer demand, have resulted in the widespread adoption of electronic transactions, record keeping, and data sharing. As one example, e-currencies, such as Bitcoin, have displaced paper currencies in millions of transactions per year. Improvements in tracking and/or knowledge management attached to such electronic transactions will continue to increase the features and options available to operators engaging in electronic transactions."}
-{"text": "1. Field of the Invention\nThe present invention relates to image-forming devices including an independent control panel which is separate from a main body of the image-forming device. In particular, the present invention relates to a copying device including an independent control panel having a display function, arranged so as to minimize the size of the copying device, and particularly to minimize the floor space required for installation of the copying device, including the main body itself and an operational area of the copying device.\n2. Description of the Related Art\nA copying device, as a known image-forming device, is shown in FIGS. 7 and 8. Other known image-forming devices include printers, facsimiles, combinations of printers and facsimiles, and the like.\nA copying device 100 shown in FIGS. 7 and 8 is provided with an automatic document-feeder 102 at an upper part of a main body 101 of the copying device 100. The automatic document-feeder 102 automatically feeds original sheets to an image reader 103, as shown in FIG. 8. The automatic document-feeder 102 includes a mounting part 102a for setting the original sheets thereon. The image reader 103 optically reads the original sheets or books that are placed upon the image reader 103. The image reader 103 is disposed at the uppermost part of the main body 101, which opposes a lower face of the automatic document-feeder 102.\nThe automatic document-feeder 102 can pivot on hinges disposed at the rear side thereof. This allows the automatic document-feeder 102 to be maintained open at any desired angle with respect to the main body 101, so that an original sheet or a book can be set on the image reader 103 for reading.\nAn original-sheet discharge-tray 102b for stacking original sheets which have been read and discharged by the automatic document-feeder 102 is provided at the right side of the upper part of the main body 101. The main body 101 is provided with a display 106a at the front side of the automatic document-feeder 102 and at the upper part of the main body 101.\nA sheet-supply tray 105 for stacking sheets for printing, when the sheets are manually supplied to the main body 101 by an operator, protrudes from the left side of the main body 101.\nWhen the operator operates the copying device 100 by using a control panel 106, the image reader 103 reads data from an original sheet or a book, an imaging unit (not shown) provided in the main body 101 transfers an image data corresponding to the read data to a sheet which was supplied from a sheet cassette 107 or the sheet-supply tray 105, and the sheet is discharged to a sheet-discharge tray 110.\nThe control panel 106 of the known copying device 100 is positioned at the front side of the upper part of the main body 101, that is, between the automatic document-feeder 102 and the operator. The control panel 106 preferably is positioned toward the front side of the copying device 100, because if it is disposed toward the rear side of the copying device 100, the control panel 106 will interfere with the easy opening and closing of the automatic document-feeder 102 or setting of the original sheets or books on the image reader 103.\nTherefore, the control panel 106 conventionally has been positioned toward the front side. Since the space provided to install the control panel 106 is limited at the front side, however, the size of the control panel 106 has had to be reduced.\nIn recent years, however, the size of the control panel 106 of known copying devices has increased. Currently, a large liquid crystal monitor including a touch panel often is provided in order to cope with the trend toward multiple functionality of the copying device and improved maneuverability thereof. As a result, an enlarged control panel 106 displaces the automatic document-feeder 102 and the image reader 103 toward the rear side of the main body 101, whereby the original sheets or books are not easily set on the image reader 103 and the size of the copying device 100 is increased.\nIn particular, it may be particularly difficult for an operator in a wheelchair, who operates the copying device 100 from a position lower than that of an operator who operates the same while standing, to set original sheets or books on the image reader 103 and handle the automatic document-feeder 102.\nMoreover, it is difficult for an operator in a wheelchair to view a control panel 106, which is substantially horizontal to the upper face of the main body 101 (the control panel 106 has poor visibility for the operator in a wheelchair).\nThe visibility of the control panel 106 may be improved by providing an operation surface of the control panel 106 that is substantially vertical, and by providing the control panel 106 toward the front side of the main body 101 of the copying device 100.\nHowever, when the control panel 106 is disposed toward the front side of the main body 101, the control panel 106 and a supporting member for the control panel 106 may interfere with the opening and closing of the automatic document-feeder 102 and with the opening and closing of a cover 104 of the main body 101 of the copying device 100. However, if the control panel 106 is disposed toward the rear side of the main body 101 of the copying device 100, remote from the operator, operability and maneuverability of the copying device 100 is reduced.\nCopying devices that attempt to overcome the above drawbacks are disclosed in, for example, Japanese Patent Laid-Open Nos. 9-6185 and 11-119498. The copying device disclosed in Japanese Patent Laid-Open No. 9-6185 is provided with a control panel 121 having a display 121a, as shown in FIG. 9, which is movable to a position appropriate for operation. The copying device disclosed in Japanese Patent Laid-Open No. 11-119498 is provided with an inclined control panel 131, as shown in FIG. 10, of which the angle of inclination is variable.\nEach of copying devices 120 and 130 shown in FIGS. 9 and 10, respectively, is provided with sheet cassettes 107 which can be drawn from the main body 101. The copying device 120 is provided with a cover 104 for opening and closing the main body 101. The automatic document-feeder 102 is provided with a mounting part 102a for setting original sheets thereon.\nThe copying devices 120 and 130 are provided with control panels 121 and 131, respectively, and displays 121a and 131a, respectively, which are positioned at the respective front sides of the automatic document feeders 102. Since the control panel 121 or 131 is positioned at the front side, the control panel 121 or 131 interferes with the opening and closing of the automatic document-feeder 102. Therefore, the display 121a of the control panel 121 must be moved or the angle of inclination of the control panel 131 must be changed each time the automatic document-feeder 102 is opened or closed.\nMoreover, it is difficult to position the automatic document-feeder 102 further toward the front side, and it is necessary to form a space around a main body 122 or 132 of the copying device 120 or 130, respectively, in which the display 121a or the control panel 131, respectively, can move. That is, there is a drawback in that the space around the copying device 120 or 130 is inefficiently used, because, when the device is operated, a larger operational space is required in addition to the space which is occupied by the copying device 120 or 130 when it is operated.\nAccordingly, it is an object of the present invention to provide an image-forming device including an independent control panel provided with a display function, in which an installation space of the device, including an operational space of a main body of the image-forming device, is minimized.\nIt is another object of the present invention to provide an image-forming device including an independent control panel, in which the control panel is easily accessible by an operator, and in particular, an operator confined to a wheelchair, without requiring excessive reach to the control panel, when the operator operates the control panel at the front side of the image-forming device.\nTo these ends, according to an aspect of the present invention, an image-forming device comprises a main body; a protrusion extending from a side surface of the main body; a support member provided either on the protrusion or on the protrusion and the main body; and a control panel including a display. The control panel is separate from the main body and supported by the support member toward a front side of the protrusion.\nThe support member of the image-forming device according to the present invention may extend above the protrusion from a rear side of the protrusion toward the front side thereof and support the control panel at the extended distal end of the support member.\nThe control panel of the image-forming device according to the present invention may be supported by the supporting member such that the control panel is capable of pivoting, rotating or turning in at least one of a vertical direction and a horizontal direction, by any arbitrary angle.\nThe protrusion of the image-forming device according to the present invention may include a control member which provides control in connection with at least one operation/function of the control panel, such as copy start, copy count, and the like.\nThe control member of the image-forming device according to the present invention may be provided at the front side of the protrusion.\nThe protrusion of the image-forming device according to the present invention may comprise a tray for stacking original sheets or sheets for printing images to be supplied to or discharged from the main body.\nThe protrusion of the image-forming device according to the present invention may comprise a mounting table for receiving thereon or therein operation manuals.\nThe control panel of the image-forming device according to the present invention is separate from the main body and supported by the support member toward the front side of the protrusion, where the support member is provided either on the protrusion or on the protrusion and the main body, whereby the maneuverability, accessibility and visibility of the control panel is improved.\nThe control panel of the image-forming device according to the present invention is supported and is disposed using a protrusion which is necessary for operation of the image-forming device. Therefore, it is not necessary to provide an additional exclusive mounting space, whereby the size of the image-forming device can be reduced.\nSince the control panel does not protrude significantly from the front face of the main body, the image-forming device can be reduced in size and the space necessary for operation can be reduced to a minimum, whereby the image-forming device can be installed in a more compact space.\nFurther objects, features and advantages of the present invention will become apparent from the following description of the preferred embodiments with reference to the attached drawings."}
-{"text": "Conventionally, for example, in automobile electric components such as automobile lamps and ECUs (Electrical Control Units), OA (office automation) equipment, household electric appliances, and medical devices, a housing containing an electric component or a control board is provided with an opening for the purpose of reducing pressure variation in the housing caused by temperature change or allowing replacement of air in the housing, and a ventilation member is attached to the opening. The ventilation member ensures ventilation between the inside and outside of the housing and also prevents foreign matters such as dust and water from entering the housing.\nFor example, Patent Literature 1 discloses a ventilation member 100 as shown in FIGS. 9A and 9B. The ventilation member 100 includes: a water-proof gas-permeable membrane 120 for covering an opening 151 of a housing 150; a support 110 supporting the water-proof gas-permeable membrane 120 and configured to be fixed to the housing 150; and a cover component 130 covering the water-proof gas-permeable membrane 120.\nThe support 110 has a tubular shape as a whole, and includes: a flange portion 112 to which the water-proof gas-permeable membrane 120 is adhered; and an insertion portion 111 projecting from the flange portion 112 and configured to be inserted into the opening 151 of the housing 150. A sealing member 140 is mounted around the root portion of the insertion portion 111. The sealing member 140 is pressed against the housing 150 by the flange portion 112, so that a gap between the ventilation member 100 and the housing 150 is sealed.\nThe cover component 130 includes: a principal wall 131 facing the water-proof gas-permeable membrane 120; and a peripheral wall 132 extending from the peripheral portion of the principal wall 131 to the vicinity of a surface of the housing 150. In addition, the principal wall 131 is provided with through holes 133 for allowing a space facing the water-proof gas-permeable membrane 120 to communicate with the outside.\nIn addition, Patent Literature 2 discloses a ventilation structure 300 in which partition walls 152 are provided on the housing 150 so as to surround a ventilation member 200 as shown in FIGS. 10A and 10B. The partition walls 152 are arc-shaped walls spaced from each other, and have a height greater than that of the ventilation member 200."}
-{"text": "Field\nEmbodiments described herein generally relate to methods and devices for self-calibration of an actuator. More specifically, embodiments described herein relate to methods and devices for self-calibration of a linear actuator configured to control a spring return valve with variable stroke.\nBackground\nCalibration of an actuator ensures proper operation and optimal performance. One approach is to use mechanical switches to detect stroke limits. Due to the mechanical contact, the accuracy, stability, and reliability of such systems may be significantly compromised. Furthermore, the mechanical switches take up additional space.\nAnother approach is to use resistors as current sensors for signal detection of motor power lines. However, due to variation in current flow, the accuracy and reliability of such sensor systems may be adversely affected. Furthermore, the sensing resistor results in poor system efficiency, resulting in using a larger actuator for the same application because the resistance must be large enough for sufficient signal sensing."}
-{"text": "One of the most CPU intensive activities associated with performing network protocol processing is the need to copy incoming network data from an initial landing point in system memory to a final destination in application memory. This copying is necessary because received network data cannot generally be moved to the final destination until the associated packets are: A) analyzed to ensure that they are free of errors, B) analyzed to determine which connection they are associated with, and C) analyzed to determine where, within a stream of data, they belong. Until recently, these steps had to be performed by the host protocol stack. With the introduction of the intelligent network interface device (as disclosed in U.S. patent application Ser. Nos. 09/464,283, 09/439,603, 09/067,544, and U.S. Provisional Application Ser. No. 60/061,809), these steps may now be performed before the packets are delivered to the host protocol stack.\nEven with such steps accomplished by an intelligent network interface device, there is another problem to be addressed to reduce or eliminate data copying, and that is obtaining the address of the destination in memory and passing that address to the network interface device. Obtaining this address is often difficult because many network applications are written in such a way that they will not provide the address of the final destination until notified that data for the connection has arrived (with the use of the \u201cselect( )\u201d routine, for example). Other attempts to obtain this address involve the modification of existing applications. One such example is the Internet Engineering Task Force (IETF) Remote DMA (RDMA) proposal, which requires that existing protocols such as NFS, CIFS, and HTTP be modified to include addressing information in the protocol headers. A solution is desired that does not require the modification of existing applications or protocols."}
-{"text": "1. Field\nA spike fixing member and a cloth treating apparatus having the same are disclosed herein.\n2. Background\nRecently, various kinds of cloth treating apparatuses other than washing machines have been developed and used. Representatively, there is a cloth treating apparatus having various functions of drying washed clothes, removing wrinkles in clothes by applying hot air or water, removing a smell remaining on the clothes, and so on. The cloth treating apparatus may be variously referred to as a cloth managing apparatus, a refresher, or a styler.\nIn Korean Patent Publication No. 10-2015-0085728, there is disclosed a cloth treating apparatus having a structure in which dried air heated while passing through a condenser by driving of a heat pump including a compressor, the condenser, an evaporator, and an expander may be supplied to a cloth accommodation space and dries clothes, and moist air in the cloth accommodation space may be dehumidified and cooled while passing through the evaporator and then circulated while passing through the condenser. Condensate water generated while the air passes through the evaporator may be collected in a container. In Korean Patent Publication No. 10-2009-0102400, there is disclosed a cloth treating apparatus in which a steam generating unit and a steam spraying unit for spraying steam into a cloth accommodation space are provided to remove wrinkles of accommodated clothes and to enhance deodorization performance.\nThe cloth treating apparatus having such a structure may be provided at an indoor space. When the cloth treating apparatus is provided at the indoor space, a horizontal control thereof is essential for a normal operation of the cloth treating apparatus. To this end, a leg which adjusts a height may be provided at a lower surface of a cabinet of the steam generating unit.\nA representative structure of the leg is disclosed in Korean Patent Publication No. 10-2015-0042948. The structure of the leg may adjust the height and may have a vibration preventing member to reduce the vibration.\nHowever, when the cloth treating apparatus is provided an indoor space on which carpet is spread, it may be difficult to stably fix the cloth treating apparatus. If it is impossible to stably fix the cloth treating apparatus, vibration may occur when the cloth treating apparatus is operated. The above references are incorporated by reference in their entirety herein where appropriate for appropriate teachings of additional or alternative details, features and/or technical background."}
-{"text": "A database-management means may comprise an optimizer engine that identifies an efficient way, known as an \u201cexecution plan,\u201d to respond to a database-access request. Such an optimizer engine may be necessary when a database-management means can derive an execution plan from a large number of possible data-access strategies, data join methods, and other implementation-specific details when in order to serve a query, generate a report, or respond to a request to update stored data. In all but the most modest configurations, selecting the most efficient execution plan may require an optimizer engine to consider a huge number of possibilities.\nMost optimizers address such tasks by analyzing a set of statistics that characterize the contents of a database. These statistics may include information about stored database objects like tables, fields, records, indices or keys, and partitions. Collecting or gathering such statistics may be a resource-intensive task that requires a distinct statistics-collection function, and this function may need to be performed frequently in order to ensure that the statistics it collects are as current as possible or, as described herein, to ensure that the statistics have not become \u201cstale.\u201d\nEach such statistics-collection or statistics-generation job may consume so many resources that it reduces database performance to unacceptable levels. When possible, these functions may be scheduled as batch jobs that run after hours in order to minimize their effect on users.\nThese conflicting requirements may create a paradox for database system designers. If a mission-critical database must be available at all times, it may be difficult to identify idle periods in which a statistics-collection task may be run with minimal effect on database availability or performance. But if statistics-collection tasks are not performed frequently enough, the statistics available to an optimizer engine may become too stale to allow the optimizer to identify execution plans that consistently ensure a efficient response to a database-access request."}
-{"text": "Understanding the reservoir fluids, hydrocarbon-generating potential, and formation characteristics of a reservoir rock formation is important to estimate the reserve and to improve the efficiency of oil field operations such as drilling, well completion, and production. To this end, various types of formation and reservoir analysis are performed using downhole tools. Further, various types of rock sample analysis are performed using laboratory tools at earth's surface. One example of formation or rock sample analysis involves nuclear magnetic resonance (NMR) phenomena.\nAn example NMR tool includes a magnet assembly that produces a static magnetic field (B0), and a coil assembly that generates radio frequency (RF) electromagnetic signals (including a perturbing magnetic field (B1)) and that detects magnetic resonance phenomena. The behavior of nuclei in presence of B0 and B1 has been correlated with formation rock and fluid properties such as the amount of hydrogen in pore space fluids and/or the porosity.\nTotal organic carbon (TOC) is an important reservoir quality indicator, particularly for unconventional source rock reservoirs. It is well-known that TOC can be obtained from laboratory mineralogy and fluid analysis using techniques such as the combination of X-ray diffraction (XRD) and X-ray fluorescence (XRF). However, laboratory tests are expensive and time consuming. Performing such tests for a large number of samples corresponding to a particular reservoir or even a single well is impractical. Nuclear logging based mineralogy analysis can also be used to estimate TOC, provided that all elements in the formation can be resolved, which is not always possible. Combining density and NMR logging measurements can also be used to estimate TOC, with an assumption that NMR logs do not include any signals from organic matters, liquid or solid, in the source rocks, an assumption not always valid.\nIt should be understood, however, that the specific embodiments given in the drawings and detailed description below do not limit the disclosure. On the contrary, they provide the foundation for one of ordinary skill to discern the alternative forms, equivalents, and other modifications that are encompassed in the scope of the appended claims."}
-{"text": "1. Field of the Invention\nThe field of this invention relates to repairing holes and rusted areas of automobile and truck body panels.\n2. Prior Art\nA major problem concerning the typical automobile whose body is fashioned from sheet metal is that certain sections of the sheet metal body tend to rust away. This rusting causes the automobile to become unsightly, which can lead to the owner's disposing of the vehicle long before the mechanical parts wear out.\nCoping with this problem has long been a challenge to those wishing to repair the automobile body.\nThe method of repairing automobile bodies has changed quite materially over a period of years. Some years ago, it was common practice to repair a dented surface or a hole in the body by flattening out the surface as well as possible by hand, and then applying molten lead to the metal surfaces after the paint had been removed. This molten lead adhered to the metal surface and was built up to a thickness sufficient to produce a smooth outer surface. The required area was sanded to provide a smooth continuous surface and the repaired area was then painted to the color of the original body.\nThis process of \"leading\" the injured areas, was not only time consuming, but it also required a relatively high degree of skill to accomplish good results. As a result, this method has fallen from general use and has been replaced by alternative methods.\nThe current methods of repairing rust damage fall into two categories which can be characterized as \"temporary\" and \"semi-permanent.\"\nThe temporary method involves applying a pressure sensitive aluminum or stainless steel tape or a fiberglass patch over the damaged area, filling over this patch with a polyester filler, priming the area and painting it. This type of repair, while quick and inexpensive, will generally last only about six (6) months. Such a repair, therefore, provides only a temporary solution to the problem which is generally less than satisfactory.\nThe semi-permanent method involves fashioning a metal patch to cover the area in question, welding or riveting the patch in place, (or in some severe cases of rust, replacing the entire section with a new section), filling over the area with polyester filler, priming and painting.\nThis method does not generally provide a fully permanent repair since it does not correct the situation which caused the rust in the first place.\nRusting in an automobile is usually caused by dirt, salt and moisture becoming trapped in a crevice or boxed-in section of the automobile body, so that the body metal is kept constantly damp and susceptible to oxidation. Thus, the rusting process occurs from the backside of the body panels and it is obvious that patches on the surface of a panel do nothing to impede the resuting process and, therefore, such patches are classified as semi-permanent.\nThis method is less than satisfactory not only because it is semi-permanent, but also because it presents the following problems.\n1. A fair amount of skill is required to fashion the metal patch so that it conforms to the original contours of the automobile body. In some instances such as the dog-leg area (the lower area where the rear fender meets the door) the compound curves are so complex that even skilled craftsmen have difficulty in consummating repairs.\n2. Skill is also required in attaching the patch or replacement section to the automobile. The repair may be riveted or welded and care must be taken that the original contours are followed. Welding will most times cause metal to distort.\n3. In most cases, the patch metal is fastened to the outside of the old metal. Thus a ridge, equal to the thickness of the patch material, is created. Since this ridge must be blended out so that the repair will be pleasing to the eye, the polyester filler must extend a relatively large distance from the repair. This necessitates that the paint be ground off a large area around the repair so that the polyester fill can be graded up to cover the patch and approximate the original contour.\n4. Welding or riveting equipment is required.\nIn many instances, repair of rusted areas is so time consuming, and the cost of skilled craftsmen is so high that the labor cost will result in a price so high that such a repair is impractical.\nClearly, a method of repairing holes and rusted areas of automobile and truck body panels which is permanent, and requires less skilled labor is greatly needed."}
-{"text": "Conventionally, numerous methods using electronic mails have been known as simple means of operating of various information processing apparatuses. For example, Japanese Patent Application Laid-Open No. 10-334002 discloses transferring an electronic mail to another network. Japanese Patent Application Laid-Open No. 6-139164 discloses updating an address book. Japanese Patent Application Laid-Open No. 6-342409 discloses changing a schedule. Further, Japanese Patent Application Laid-Open No. 5-233493 discloses management of computer.\nHowever, in the above respective conventional technologies, application of such technologies to an information processing apparatus capable of voice output is not disclosed.\nOn the other hand, it is generally possible to voice-output a transmitted electronic mail on the receiving apparatus. However, the voice output cannot be performed at a sender's intended timing. Rather, the timing of the voice output of the transmitted electronic mail is determined by the receiving apparatus."}
-{"text": "The present invention relates to a multi-piece golf ball having at least a three-layer construction composed of a core, an intermediate layer and a cover layer. More specifically, the invention relates to a multi-piece golf ball having a high rebound resilience and durability, and having an excellent feel on impact. The invention also relates to a multi-piece golf ball which has a low spin rate and can travel a good distance even when hit at a low head speed.\nThe golfing population has grown in recent years, and the desires of players with regard to golf balls have become increasingly diverse and personalized. To address such desires, various investigations are being conducted on ball construction. Most players are amateur golfers with head speeds (HS) in a range of generally 30 to 45 m/s. Amateur golfers derive greater enjoyment from playing the game when the ball travels farther on shots with a driver and has a good, soft feel at the moment of impact. In addition, because amateur golfers have a head speed that is somewhat lower than that of professionals, it tends to be harder for them to get enough lift on the ball. Amateurs thus use clubs having a relatively high loft angle, as a result of which balls hit with a number one wood (W#1) tend to take on too much spin and acquire an excessive spin rate.\nJP-A 2005-218858 teaches a golf ball having a core, an intermediate layer that encloses the core, and a cover that encloses the intermediate layer. By having the intermediate layer and the cover each formed of an ionomer resin-containing thermoplastic resin, and by concurrently optimizing the respective values for the core diameter, deflection of the core under 100 kg of loading, intermediate layer thickness, Shore D hardness of intermediate layer, cover thickness, Shore D hardness of cover, deflection of the golf ball under 100 kg of loading, and the difference obtained by subtracting the deflection of the golf ball under 100 kg of loading from the deflection of the core under 100 kg of loading, such a golf ball achieves a combination of a good, high rebound, low spin and good feel when used by the ordinary amateur golfer who values distance and feel more than spin performance.\nJP-A 2004-16583 discloses a golf ball having, in order from the inside: a solid core of at least one layer, an intermediate layer, and a cover. In this ball, a good rebound and a soft feel on impact are achieved by optimizing the respective hardnesses and thicknesses of the core, intermediate layer and cover.\nHowever, these golf balls leave room for further improvement in their distance of travel and feel on impact. In particular, the challenge has been how to lower the spin rate on the ball when hit with a driver in the low head speed range typical of the amateur golfer, and thereby improve the distance traveled by the ball."}
-{"text": "1. Field\nExample embodiments relate to an automatic packing apparatus which automatically packs products in boxes and an automatic packing method using the same.\n2. Description\nIn general, after various products are manufactured, the products are packed in boxes in order to prevent the products from being damaged and to ship/handle the products.\nThe packing boxes are made of a material, such as thick corrugated cardboard, so as to protect the products from external impact. The manufactured packing boxes are stored in an unassembled state in which the boxes are unfolded in a flat shape so as to reduce a space occupied thereby before being used to pack the products. Then, when the packing boxes are used to pack the products, the packing boxes are assembled (folded) to pack the products.\nA worker folds a packing box into an original shape, and then covers a product with the packing box, thereby packing the product in the packing box. Here, the product is covered with the box downwardly from above. Packing of a large-sized home appliance product, such as a refrigerator, a TV or an air conditioner, in such a box requires at least two workers.\nAs described above, workers manually pack products in boxes and thus packing requires considerable manpower, and a large amount of time is consumed packing each product and thus productivity is reduced."}
-{"text": "1. Field of the Invention\nThe present invention relates to a plasma processing apparatus, and more specifically, relates to a plasma processing apparatus capable of suppressing the occurrence of particles caused by reaction products.\n2. Description of the Related Art\nIn the field of manufacturing semiconductor devices, plasma processing apparatuses using plasma are adopted as means for realizing microfabrication to respond to the increasing demands for highly precise high-speed microfabrications (such as dry etching), to realize enhanced performance of semiconductor integrated circuit elements and cut down costs. Among such arts, the art of dry etching is often applied, in which processing gas is used to process films deposited on the sample and evacuate the residual gas in the processing chamber via a turbo-molecular pump or the like to the exterior. Various kinds of gases (such as Ar, Cl2 and HBr) can be used as the processing gases, but some gases are not sufficiently evaporated, depending on the kind of films deposited on the sample being processed. In such cases, gases are not sufficiently discharged and become the cause of reaction products, which deposit on the inner wall of the vacuum processing chamber and on the inner wall of an inductive vacuum processing chamber lid disposed on the upper portion of the vacuum processing chamber. When excessive reaction products are deposited on the inner wall of the vacuum processing chamber and the inner wall of the vacuum processing chamber lid, deposition films may fall off from the surface of the inner wall and attach to the surface of the sample as particles.\nOne type of prior art plasma processing apparatuses is an inductive plasma processing apparatus having a high frequency power supply connected to a coil-like induction antenna disposed on an outer circumference of a vacuum processing chamber lid for supplying high frequency power to the coil-like induction antenna, to thereby generate plasma. The drawback according to this type of plasma processing apparatus adopting an induction antenna is that the bonding state of the induction antenna and the plasma within the vacuum processing chamber lid is varied by the reaction products deposited on the inner wall of the vacuum processing chamber and the vacuum processing chamber lid, by which the etching rate, the uniformity of etching, the perpendicularity of the etching profile and the state of deposition of reaction products on the etched side walls is varied with time.\nNext, Japanese patent application laid-open publication No. 2004-235545 (patent document 1) discloses the following method as the method for preventing reaction products from depositing on the inner wall of the vacuum processing chamber lid. A Faraday shield capacitively coupled with plasma is arranged between induction antennas disposed at the outer circumference of a vacuum processing chamber lid and plasma, and high frequency power is supplied to the Faraday shield to generate self bias to the inner wall of the vacuum processing chamber lid, drawing ions in the plasma by the electric field within the ion sheath and sputtering the surface of the inner wall of the vacuum processing chamber lid. The sputtering is used to suppress or remove the deposition of reaction products on the inner wall of the vacuum processing chamber lid, and perform cleaning of the inner wall of the vacuum processing chamber lid. In the present specification, the vacuum processing chamber lid is airtightly fixed to the upper portion of the vacuum processing chamber.\nAs described, a self bias is generated on the surface of the inner wall of the vacuum processing chamber lid via the high frequency voltage applied to the Faraday shield, by which the electric field of the ion sheath is increased and the energy of ions reaching the inner wall of the vacuum processing chamber lid is enhanced. Therefore, ion sputtering via ions becomes prominent, enabling more reaction products to be cleaned. Further, most appropriate cleaning becomes possible by controlling the high frequency voltage supplied to the Faraday shield.\nFurther, the following method is disclosed in Japanese patent application laid-open publication No. 2005-259836 (patent document 2) as a method for cleaning the whole surface of the inner wall of the vacuum processing chamber lid in a uniform manner. Patent document 2 discloses a plasma etching apparatus having a Faraday shield divided into multiple parts and providing independent high frequency power supplies respectively to each of the divided Faraday shields, to thereby enable control of the voltages applied to the respective Faraday shields.\nAccording to the disclosed plasma processing apparatus, in order to remove the reaction products generated from samples during etching and thickly deposited on a center portion of the vacuum processing chamber lid (top plate) closest to the sample surface, a high frequency power which is higher than the power supplied to the outer circumference portion of the top plate is supplied to the divided Faraday shield disposed at the center portion of the top plate. According to this arrangement, reaction products deposited both on the center portion and the outer circumference portion of the top plate can be cleaned appropriately.\nFIG. 10 shows the state of deposition of reaction products attached to the inner wall of the vacuum processing chamber lid in a plasma processing apparatus. When etching is performed, a large amount of reaction products is generated from the sample being processed, and the generated reaction products are attached to and deposited on the inner wall of the vacuum processing chamber lid (hereinafter referred to as a window 12) and the inner wall surface of a vacuum processing chamber 2 as a film (18).\nIn contrast, at an outermost circumference portion (portion B) of the inner wall surface of the window 12, the self bias is reduced since the vacuum processing chamber 2 is grounded, so that the energy of ions with respect to portion B is reduced. Ion sputtering rarely occurs in portion B, and therefore, reaction products remain in portion B. As a result, the inner wall surface of the window 12 is divided into a region (portion A) where the reaction products are completely removed, and a region (portion B) where reaction products easily remain.\nIn order to eliminate residual reaction products in portion B, it is necessary to set the high frequency voltage applied to the Faraday shield 17 (Faraday shield voltage: hereinafter abbreviated as FSV) as high as possible to remove the reaction products deposited on the inner wall of the window 12. According to this method, however, the inner wall surface of the window 12 of portion A is chipped excessively.\nAs described, if the inner surface of the window 12 is chipped and the speed of consumption of the window 12 itself is increased, the replacement cycle of the window 12 is shortened. Further, the chipped window 12 becomes the source of reaction products and re-deposits on the wafer surface, causing instability of etching rate of the process performance and in-plane unevenness of the wafer etching rate and the like. Another drawback is that the chipped surface of the window 12 turns into particles and attach to the sample surface. Therefore, in order to overcome the above-mentioned problems, it is necessary to remove the reaction products deposited on the whole inner wall surface of the window 12 in a most suitable manner.\nPrior arts such as that disclosed in patent document 2 enables to remove the reaction products deposited on the whole inner wall surface of the window 12 in a suitable manner, but since independent high frequency power supplies must be disposed in response to the divided Faraday shields and antennas, the application thereof to industrially-applied plasma processing apparatuses causes increase of costs and requires a large mounting space.\nFurthermore, there is a drawback that cannot be solved by the prior art disclosed in patent document 2. This drawback will be described with reference to FIG. 5. Since the present drawback is not related to the division number of Faraday shields, the plasma processing apparatus illustrated in FIG. 5 has a Faraday shield that is not divided into multiple parts.\nIn FIG. 5, the dotted lines with arrows show high frequency currents flowing from a Faraday shield 17 into the vacuum processing chamber 2 during plasma processing. The Faraday shield 17 is capacitively coupled with plasma, and the high frequency current from the Faraday shield 17 flows through the plasma as shown by the arrows into the vacuum processing chamber 2 being grounded.\nNow, when the high frequency voltage applied to the Faraday shield 17 is denoted by Va, the area of the Faraday shield 17 is denoted by Sa, the high frequency voltage generated in the ion sheath on the surface of the inner wall of the vacuum processing chamber 2 is denoted by Vb, and the surface area of the ion sheath is denoted by Sb, the following expression is satisfied, as taught in cited document (Michael A. Lieberman, Allan J. Lichetenberg, ED Research, \u201cPrinciples of Plasma Discharges and Materials Processing\u201d Chapter 11).\n                                          V            b                                V            a                          =                              (                                          S                a                                            S                b                                      )                    n                                    [                  Expression          \u2062                                          \u2062          1                ]            \nExpression 1 shows that a self bias is not only generated at the inner wall surface of the window 12 by the high frequency voltage applied to the Faraday shield 17, but also generated at the inner wall surface of the vacuum processing chamber 2 into which high frequency current flows.\nTherefore, the inner wall of the vacuum processing chamber 2 is sputtered via the ions accelerated by Vb. As described, when the inner wall of the vacuum processing chamber 2 is sputtered, the lifetime of the respective parts composing the interior of the vacuum processing chamber 2 is reduced, and the frequency of replacement of the parts is increased. Thus, the running costs of the plasma processing apparatus cased by replacement of parts is increased. Even further, the chipped inner wall of the vacuum processing chamber 2 becomes the source of particles, by which the mass production performance of the plasma processing apparatus is deteriorated."}
-{"text": "The present invention is in the field of preparing self-assembled monolayers on a metal using electrolysis.\nMolecular self-assembly of monolayers (SAMs) provides a simple method to form highly ordered two-dimensional organic assemblies. Among the various systems that display this behavior, SAMs formed by the chemisorption of alkanethiols on gold to produce strong gold-thiolate bonds is particularly convenient because of the ease of their preparation. [1-3] Well-ordered SAMs form spontaneously on gold surfaces within a short period upon immersion of the surface in either a dilute solution or the vapor of an alkanethiol of interest. These monolayers have been used as model systems for fundamental studies of wettability [4], adhesion [5], biocompatibility [6], fouling [7], as well as serving as the basis for building multilayers [8] and bio- and analytical sensors [9], for the immobilization of biomolecules [10], and for preparing patterned surfaces [11].\nOne of the most attractive characteristics of self-assembled monolayers (SAMs) is the facility of their preparationxe2x80x94simply exposing a substrate surface to an appropriate adsorbate-[1-3] though this facility prevents selective formation at particular surfaces in the presence of others of the same composition. While contact printing of SAMs has proven useful in preparing mesoscale patterns on various substrates [12], formation of monolayers on only particular features of a pre-existing patternxe2x80x94such as an electrode array, an integrated circuit, or a MEMS devicexe2x80x94has remained an elusive synthetic challenge. Hence, treating an array of identical gold electrodes with a solution of an alkanethiol or alkyl disulfide would result in formation of a monolayer on all of the electrodes. One approach to controlling self-assembly in these systems has focused on the inability of thiolate ions to covalently bind the gold surface directly, without concomitant oxidation. [13] Hence, gold electrodes immersed in a solution of thiolate ions will only adsorb a SAM if held at a sufficiently high potential. Electrochemistry, in this case reduction, can also be used to remove SAMs that had previously been adsorbed. [14] Hence, an alternative strategy for producing patterns involves allowing indiscriminate self-assembly to occur, followed by articulation of a pattern by the selective removal of the SAM from certain substrates.\nAlkylthiosulfates, also known as Bunte salts, can be used to synthesize disulfides by oxidation [15], acidic hydrolysis [16-17], or alkaline degradation [18]. Disulfides also can be formed from Bunte salts electrochemically. [19-20] This method has been extended to form polydisulfides from xe2x80x9cdoublexe2x80x9d Bunte salts, molecules carrying two thiosulfate groups, using electrochemistry with gold electrodes. [21]\nDisclosed is a method of preparing self-assembled monolayers on a metal comprising electrolyzing a thiosulfate compound in a solvent, where the electrolysis is performed at a voltage for a period of time.\nAlso disclosed is a method of preparing self-assembled organic monolayers on a metal comprising (a) contacting said metal with a solution comprising an organic thiosulfate compound, and (b) electrolyzing said organic thiosulfate compound by applying on said metal sufficiently high anodic potential for sufficient time to result in the oxidative self-assembly of said monolayers on said metal.\nFurther, a method is disclosed for the selective formation of self-assembled organic monolayers on a first metal electrode in the vicinity of a second metal electrode, comprising (a) contacting said metal electrodes with a solution comprising an organic thiosulfate compound under conditions such that chemisorption of said organic thiosulfate compound onto said first and second electrodes does not occur, and (b) electolyzing said organic thiosulfate compound by selectively applying on said first metal electrode sufficiently high anodic potential for sufficient time to result in the oxidative self-assembly of said organic monolayers on said first electrode."}
-{"text": "One of the known problems associated with many industrial factories operating under elevated temperatures is the generation of excess heat emitted from the facility, which is not effectively utilized.\nAlthough there could be a substantial potential of energy content in the outgoing gases, still, the typical problem associated with utilizing this potential is the fact that the hot gaseous feed stream is a by-product generated on a transient basis, e.g., an outcome of batch processes, non-constant generation of the by-product hot gases, and the like. Thus, it would be virtually impossible to rely on an essentially non-constant hot gases feed or even intermittent feed of hot gases, for running a continuous process using this heat, which is desired from economical point of view.\nOne way of solving this problem, is by including a storage in the system for storing the hot gases leaving the process and drawing a constant hot gaseous feed from that storage. However, this solution is impractical due to the costs associated with the need to store and compress the gases at elevated temperatures.\nAnother way to obtain an intermediate storage is to use a phase changing material in which the heat can be stored. However, once again the costs involved with such a process and the loss of process efficiency (due to the reduction of the temperature in the storage and retrieval processes), render this solution impractical.\nYet another way to obtain such a storage is by using materials which have high heat capacity (e.g. refractory bricks through which hot gas carrying the excess heat flows) and by heating this media, the excess heat is stored as sensible heat. Again, this solution suffers from similar disadvantages as the solutions referred to above.\nIn view of that, there is a need to obtain a solution that would enable utilizing the excess heat and consequently to increase the overall efficiency of the process in which the excess heat is generated."}
-{"text": "Concrete slabs on grade, and other concrete structures and objects, traditionally are made by building a form or a mold. The forms and molds are usually made from wood, plywood, metal and other structural members. Unhardened (plastic) concrete is placed into the space defined by opposed spaced form members. Once the concrete hardens sufficiently, although not completely, the forms are removed leaving a concrete wall or other concrete structure, structural member or concrete object exposed to ambient temperatures. For some applications, such as runways, roads and highways, slip forms are used to continuously place long lengths of concrete. For both conventional forms, molds and slip forms, the unprotected concrete slabs, structures or objects are exposed to the elements during the remainder of the concrete curing process after the forming process. The exposure of the concrete to the elements, especially temperature variations, often makes the curing of the concrete a slow process and the ultimate strength difficult to control or predict. To compensate for these losses, larger amounts of portland cement are used than otherwise would be necessary. Larger cement fractions, particularly for high surface area-to-volume ratio elements, such as these, can increase the likelihood of shrinkage, cracking, and curling.\nThe curing of plastic concrete requires two elements, water and heat, to fully hydrate the cementitious material. The curing of plastic concrete is an exothermic process. This heat is produced by the hydration of the portland cement, or other hydraulic cements, that make up the concrete. Initially, the hydration process produces a relatively large amount of heat. As the hydration process proceeds, less hydration occurs thereby resulting in the production of less heat. At the same time, moisture in the concrete is lost to the environment. If one monitors the temperature of concrete during the curing process, it produces a relatively large increase in temperature which then decreases over time. This chemical reaction is temperature dependent. That is, the hydration process, and consequently the strength gain, proceeds faster at higher temperature and slower at lower temperature. During conventional concrete curing, first, the heat is lost; then, moisture loss makes it difficult for the cementitious material to fully hydrate, and, therefore, the concrete may not achieve its maximum strength.\nConcrete in conventional concrete forms or molds is typically exposed to the elements. Conventional forms or molds provide little insulation to the concrete contained therein. Concrete placed in forms for slabs on grade, roads, runways, and the like, is usually placed on a sheet of polyethylene placed on the ground first. Therefore, the bottom of the concrete slab is in contact with the ground, which absorbs the heat of hydration from the concrete. The top of the concrete slabs are sometimes covered with plastic sheathing to prevent some of the moisture loss to the environment. This also allows the concrete to lose the heat of hydration from the top of the slab surface to the air. Although during winter an insulated blanket may be placed on top of the concrete slab, road or runway to keep the concrete from freezing, heat produced within the concrete form or mold due to the hydration process usually is lost to the ground or to the air through a conventional concrete form or mold relatively quickly. Thus, the temperature of the plastic concrete may initially rise 20 to 40\u00b0 C. or more above ambient temperature due to the initial hydration process and then fall relatively quickly to ambient temperature, such as within 12 to 36 hours. This initial relatively large temperature drop may result is concrete shrinkage, curling and/or concrete cracking. The remainder of the curing process is then conducted at approximately ambient temperatures, because the relatively small amount of additional heat produced by the remaining hydration process is relatively quickly lost to the ground or to the air through the uninsulated concrete form or mold. The concrete is therefore subjected to the hourly or daily fluctuations of ambient temperature from hour-to-hour, from day-to-night and from day-to-day. Curing concrete under ambient temperature conditions is not as significant a problem during summer temperature. It is cool or cold weather conditions that cause the most significant trouble for properly curing concrete.\nFailure to cure concrete under ideal temperature and moisture conditions affects the ultimate strength and durability of the concrete. In colder weather, concrete work may even come to a halt since concrete will freeze, or not gain much strength at all, at relatively low temperatures. By definition (ACI 306), cold weather conditions exist when \u201c . . . for more than 3 consecutive days, the average daily temperature is less than 40 degrees Fahrenheit and the air temperature is not greater than 50 degrees Fahrenheit for more than one-half of any 24 hour period.\u201d Therefore, in order for hydration to take place, the temperature of concrete must be above 40\u00b0 F.; below 40\u00b0 F., the hydration process slows and at some point may stop altogether. It is typically recommended that concrete by moisture cured for 28 days to fully hydrate the concrete. However, this is seldom possible to achieve in commercial practice.\nInsulated concrete form systems are known in the prior art and typically are made from a plurality of modular form members. U.S. Pat. Nos. 5,497,592; 5,809,725; 6,668,503; 6,898,912 and 7,124,547 (the disclosures of which are all incorporated herein by reference) are exemplary of prior art modular insulated concrete form systems. Full-height insulated concrete forms are also known in the prior art. U.S. Patent Application Publication No. 2011/0239566 and applicant's co-pending patent application Ser. No. 13/247,133 filed Sep. 28, 2011 (the disclosure of which are both incorporated herein by reference in their entirety) discloses a full-height vertical insulated concrete form. Insulated concrete forms reduce the heat transmission to and from the concrete within such forms. However, vertical insulated concrete forms are not useful for forming slabs on grade, such as runways, road and highways and previously have not been proposed to be used in such applications.\nConcrete insulating blankets are known in the art. Electrically heated insulating blankets are also known in the prior art, such as those disclosed in U.S. Pat. Nos. 7,183,524 and 7,230,213. Such concrete insulating blankets and electrically heated insulating blankets are known for use in northern climates for thawing frozen ground and/or preventing curing concrete from freezing. It is known that plastic concrete will not cure satisfactorily at temperature below 50\u00b0 F. However, such electrically heated blankets are designed to provide a constant amount of heat to the plastic concrete and are used only for the purpose of preventing the concrete from freezing in cold weather.\nTherefore, it would be desirable to produce a concrete forming or molding system for slabs on grade, such as runways, road and highways. It would also be desirable to provide a concrete curing system can be used specifically to cure slabs on grade, such as runways, roads and highways. It would also be desirable to provide a concrete curing system that accelerates concrete maturity or equivalent age to achieve improved concrete strength, particularly early concrete strength. Furthermore, it would be desirable to provide a concrete curing system for slabs on grade, such as runways, roads and highways, so that the concrete cures more quickly, is less permeable, more flexible, stronger, more durable and less prone to cracking and curling. It would also be desirable to produce concrete slabs on grade, such as runways, roads and highways, that are more environmentally friendly."}
-{"text": "In the supply of ink into the nozzle head in an ink jet recorder, it is necessary that pressure of the ink be maintained constant and air bubbles be not mixed in the ink. As a matter of fact, however, it is said that the concentration of air (oxygen) dissolved in the ink in the state that the ink is only exposed to the air is increased to 5 to 6 ppm from 1 ppm for a period of 30 minutes. And, the entrained air accumulates at such as a bent area in the piping to become air bubbles which in the nozzle head cause an undischarge of ink or a fluctuation in ink pressure which in turn brings about a poor or defective printing.\nAccordingly, the ink furnishing unit of this type in the prior art includes a deaerator that is provided for a recovery tank which reserves the ink recovered from the purge section and abounding with air bubble (see, for example, JP P2003-127435 A) or a deaerator that is disposed midway in the ink supply circuit (see, for example, JP P2004-17517 A).\nIn the technique shown in JP P2003-127435 A and designed to permit recycling the ink recovered at the purge section, it is difficult to deaerate the entire recovery tank uniformly so that it has not been possible to remove air bubbles completely from the ink in the recovery tank.\nWhile in this regard the technique shown in JP P2004-17517 A permits the entire volume of ink furnished from the supply tank into the main tank to be evenly deaerated only by passing through its piping, this techniques does not allow such ink to be directly furnished to the nozzle head. Since the ink is once stored in the air-open main tank for stabilization of its supply pressure and then is furnished to the nozzle head, the amount of air entrained in the ink while in the main tank increases with the lapse of time, causing the problem that the ink discharge at the nozzle head tends to be defective.\nWith these taken into account, it is an object of the present invention to provide an ink furnishing unit in an ink jet recorder that can efficiently deaerate the ink furnished to the nozzle head including the ink recovered from the purge section so that the defective printing due to the mixing of air bubbles in the ink may be prevented."}
-{"text": "Normally, in continuous casting plants, the molten metal produced, for example from an electric furnace, or from a converter, is subjected to a series of pouring passages between containers, prior to being cast in the ingot mould to form ingots, slabs, bars, strips and others.\nA typical example is given by the application of the technology of continuous strip casting in which the steel coming from the production oven is collected in a ladle, from which it is unloaded into one or more tundishes, each of which feeds one or more melts above two cooled, counter rotating crystallising rolls, which constitute the ingot mould in which the metal solidifies and exits in the form of finished product, for example a strip, a bar or otherwise. Between the tundish and the ingot mould there can also be provided an under tundish, or another discharging device, from which the liquid metal is dosed into the ingot mould.\nIt is known that the quality of the final product, and also the execution of the process itself can be compromised by the contact of air with the molten metal. For example, the oxygen in the air can combine with elements dissolved in the steel forming inclusions, which deteriorate the quality of the steel, as also with dissolved oxygen itself. In case the steel is used in continuous casting, for example in a dual roller plant, the oxygen combined with the iron can produce scale which deposits on the rollers, locally altering, amongst others, the heat exchange, with serious consequences for the final product. Nitrogen can also provoke the formation of precipitates which compromise the quality of the product.\nIn the containers into which the molten metal passes, there is generally maintained a protective atmosphere, generally through the feeding of an inert gas, for example argon.\nThe various passages of the metal from one container to another are critical, with regard to the problem of possible contamination with air, bearing also in mind that the various containers must be emptied in between, both for the normal movement during the metal production operation, as for example in the case of the tundishes, and to allow for the substitution of the parts in refractory material, such as, for example the discharging sleeves or tubes through which the molten metal flows, subject to rapid wear and corrosion.\nDischarging devices in the known art are not easy to operate and require delicate operations to allow the approaching of two containers for joining. Sometimes the impacts produced during the coupling damage the parts of the discharging device. Therefore, the need is felt of providing discharging devices, to interpose between the various containers in which the molten metal passes, avoiding contact between the air and the metal and allowing a rapid and precise coupling between the containers, in addition to a likewise rapid uncoupling."}
-{"text": "The present invention is related to a rock climbing apparatus, and more particularly to a lightweight rock climbing apparatus which has sufficient strength and can be easily assembled for personal or domestic rock climbing exercise.\nRock climbing has become a popular activity. However, natural outdoor rock field can be hardly found. Moreover, there is safety problem in real rock climbing. Therefore, many artificial rock fields have been built. However, the artificial rock field necessitates larger space and the equipments of the artificial rock field are expensive. Therefore, the cost for such rock climbing exercise is high.\nU.S. Pat. Nos. 5,549,195, 5,125,877 and 5,919,117 disclose rock climbing apparatuses. In each of those rock climbing apparatuses, a caterpillar structure is arranged on a frame body. The caterpillar structure is connected with a transmission mechanism for restricting the rotational speed of the caterpillar structure. Two sides of the caterpillar structure are respectively disposed in the rails of the frame body via chains or rollers, whereby the caterpillar structure can circularly revolve along the rails. Multiple rocks are arranged on the caterpillar structure for a user to climb for rock climbing exercise.\nIn the caterpillar structure of the above rock climbing apparatus, simply two sides of the caterpillar structure bear the force, while the middle of the caterpillar structure totally lacks support structure. The rocks are positioned on the middle of the caterpillar structure. Therefore, when a user climbs the rocks, the application force will totally fall onto the middle of the caterpillar structure. In order to avoid swing or deformation of the caterpillar structure, the caterpillar structure is designed with quite strong structure for bearing all the application force in climbing.\nFor example, U.S. Pat. No. 5,549,195 discloses a caterpillar structure composed of multiple solid planks. U.S. Pat. No. 5,125,877 discloses a caterpillar structure composed of multiple metal frames pivotally connected with each other. A board body is inlaid in each metal frame. U.S. Pat. No. 5,919,117 discloses a caterpillar structure composed of multiple extruded board bodies pivotally connected with each other. In order to avoid deformation of the caterpillar structures, in all the above rock climbing apparatuses, the caterpillar structures are designed with quite strong and complicated structures. This results in that the rock climbing apparatus has a weight up to several hundred kilograms. Therefore, it is hard to move the rock climbing apparatus. In other words, such rock climbing apparatus is unsuitable for personal or domestic rock climbing exercise.\nFurthermore, the frame body of the rock climbing apparatus is a huge and heavy structure for supporting the heavy caterpillar. Therefore, the rock climbing apparatus can be hardly assembled or transferred.\nBesides, the caterpillar structure has gaps between the board bodies and the frames. These gaps are enlarged when the caterpillar is curved. A user's hairs, clothes, trousers or shoelaces may be chucked in the gaps to result in accident.\nU.S. Pat. No. 6,860,836 of this applicant discloses a relatively lightweight rock climbing apparatus. A circularly rotating mechanism is mounted between two frames of a bed. The circularly rotating mechanism includes several sprockets. Each chain has multiple lugs arranged at intervals for mounting multiple transverse beams in parallel to the shaft. Two ends of each transverse beam are respectively inserted in open sides of the two frames. The transverse beams connected with the chains can be circularly revolved along the inner edges of the open sides of the frames. The chains and the transverse beams are wrapped with serially connected soft pads for shading the chains and the transverse beams. Multiple rocks are locked on outer sides of the soft pads corresponding to the lugs of the chains. Several U-shaped stopper boards are fixed between the frames corresponding to the chains, whereby the chains can lean on the stopper boards.\nThe chains can circularly revolve in a fixed path between the frames. The rocks are arranged on the chains which serve as the main supports for the rocks. Multiple transverse beams are arranged on the chains at intervals for bearing the outward pulling force exerted onto the rocks. The chains lean on the U-shaped seats for bearing the inward pushing force exerted onto the rocks. The above structure is simplified and the weight thereof is minified.\nHowever, in practice, the above rock climbing apparatus still has some shortcomings. The rocks are fixed on the chains. When the rocks are forced inward, the lateral sides of two links of the chain abut against the U-shaped seat. When the chain revolves, the links and the U-shaped seat will abrade each other to produce noise. After a period of use, the links and the U-shaped seat will be worn out. This will affect the smoothness of the revolution and even the strength of the chain. In addition, when the rocks are forced outward, two ends of the transverse beam will abrade the open sides of the frames. Similarly, the transverse beam and the inner sides of the frames will be worn out and noise is produced.\nMoreover, the above rock climbing apparatus is lightweight and suitable for personal or domestic use. However, the frames still have a considerable length exceeding the height of an adult for a user climb. The frames are divided into two segments. However, the segments can be hardly conveniently DIY assembled. Furthermore, there is no support structure between the frames. The entire structure is simply supported by the bed. Therefore, the strength of the rock climbing apparatus is still insufficient.\nBesides, the transverse beams are arranged at intervals. Therefore, a gap is defined between two transverse beams. Although the transverse beams are wrapped with the soft pads, a user still often touches the gaps when climbing. At this time, the soft pads will be depressed to contact with the chains or the skeleton inside the frame. Under such circumstance, the inner faces of the soft pads will be abraded and worn out by the chains or the skeleton of the frame.\nIn addition, the soft pads must have a certain tightness to keep the surface smooth. However, the transverse beam is wrapped with the soft pad and fixed on the chain via the rock. When the transverse beam is revolved to the sprocket, the transverse beam will be tangential to the circumference of the sprocket. Accordingly, two ends of the transverse beam will lift the soft pad as shown in FIG. 12. At this time, the soft pad will be tensioned. In some cases, the soft pad will be broken in the position where the rock is fixed. In some serious cases, the chain will be unable to revolve. In the case that the soft pad is made of a resilient material, the inner face of the soft pad will be abraded by the chain or the skeleton of the frame. As a result, after a period of use, the soft pad will elastically fail due to the continuous tension."}
-{"text": "In order to quickly review and analyze a video, for example a movie, a recorded sporting event or a news broadcast, a summary of the video can be generated. A number of techniques are known for summarizing uncompressed and compressed videos.\nThe conventional practice is to first segment the video into scenes or \u2018shots\u2019, and then to extract low and high level features. The low level features are usually based on syntactic characteristics such as color, motion, and audio components, while the high level features capture semantic information.\nThe features are then classified, and the shots can be further segmented according to the classified features. The segments can be converted to short image sequences, for example, one or two seconds \u2018clips\u2019 or \u2018still\u2019 frames, and labeled and indexed. Thus, the reviewer can quickly scan the summary to select portions of the video to playback in detail. Obviously, the problem with such summaries is that the playback can only be based on the features and classifications used to generate the summary.\nIn order to further assist the review, the segments can be subjectively rank ordered according to a relative importance. Thus, important events in the video, such as climactic scenes, or goal scoring opportunities can be quickly identified, see, Fujiwara et al. \u201cAbstractive Description of Video Using Summary DS,\u201d Point-illustrated Broadband+Mobile Standard MPEG Textbook, ASCII Corp., p. 177 FIGS. 5-24 Feb. 11, 2003, also \u201cISO/IEC 15938-5:2002 Information technology\u2014Multimedia content description interface\u2014Part 5: Multimedia Description Schemes,\u201d 2002. After an important video segment has been located, the viewer can use fast-forward or fast-reverse capabilities of the playback device to view segments of interest, see \u201cDVR-7000 Instruction Manual,\u201d Pioneer Co., Ltd., p. 49, 2001.\nAnother technique for summarizing a news video uses motion activity descriptors, see U.S. patent application Ser. No. 09/845,009, titled \u201cMethod for Summarizing a Video Using Motion Descriptors,\u201d filed by Divakaran, et al., on Apr. 27, 2001. A technique for generating soccer highlights uses a combination of video and audio features, see U.S. patent application Ser. No. 10/046,790, titled \u201cSummarizing Videos Using Motion Activity Descriptors Correlated with Audio Features,\u201d filed by Cabasson, et al., on Jan. 15, 2002. Audio and video features can also be used to generate highlights for news, soccer, baseball and golf videos, see U.S. patent application Ser. No. 10/374,017, titled \u201cMethod and System for Extracting Sports Highlights from Audio Signals,\u201d filed by Xiong, et al., on Feb. 25, 2003. Those techniques extract key segments of notable events from the video, such a scoring opportunity or an introduction to a news story. The original video is thus represented by an abstract that includes the extracted key segments. The key segments can provide entry points into the original content and thus allow flexible and convenient navigation.\nThere are a number of problems with prior art video recording, summarization and playback. First, the summary is based on some preconceived notion of the extracted features, classifications, and importance, instead of those of the viewer. Second, if importance levels are used, the importance levels are usually quantized to a very small number of levels, for example, five or less. More often, only two levels are used, i.e., the interesting segments that are retained, and the rest of the video that is discarded.\nIn particular, the hierarchical description proposed in the MPEG-7 standard is very cumbersome if a fine quantization of the importance is used because the number of levels in the hierarchy becomes very large, which in turn requires management of too many levels.\nThe MPEG-7 description requires editing of the metadata whenever the content is edited. For example, if a segment is cut out of the original content, all the levels affected by the cut need to be modified. That can get cumbersome quickly as the number of editing operations increases.\nThe importance levels are highly subjective, and highly context dependent. That is, the importance levels for sports videos depend on the particular sports genre, and are totally inapplicable to movies and news programs. Further, the viewer has no control over the length of the summary to be generated.\nThe small number of subjective levels used by the prior art techniques make it practically impossible for the viewer to edit and combine several different videos based on the summaries to generate a derivate video that reflects the interests of the viewer.\nTherefore, there is a need to record and reproduce a video in a manner that can be controlled by the viewer. Furthermore, there is a need for specifying importance levels that are content independent, and not subjective. In addition, there is a need to provide more than a small number of discrete importance levels. Lastly, there is a need to enable the viewer to generate a summary of any length, depending on a viewer-selected level of importance."}
-{"text": "1. Field of the Invention\nThe present invention relates to a display device and a method of manufacturing the same.\n2. Description of the Related Art\nA micro-electro-mechanical system (MEMS) display is a display expected to replace a liquid crystal display (see Japanese Patent Application Laid-open No. 2008-197668). This display differs from a liquid crystal shutter type display utilizing polarization, and displays an image by opening and closing a light transmissive window using a mechanical shutter system. A shutter is formed of a thin film. Vertical and horizontal sizes of one shutter forming one pixel are in the order of several hundred micrometers, and a thickness thereof is in the order of several micrometers. One shutter is opened/closed to enable ON/OFF operation for one pixel. The shutter is operated by an electrostatic attractive force.\nThe shutter is arranged in a space surrounded by a sealing member between a pair of light transmissive substrates, and the space is filled with oil. The oil is used to prevent a spring for driving the shutter from sticking, and to reduce a difference in optical refraction index with respect to the light transmissive substrates.\nThe oil is injected from an injection port corresponding to an opening of the sealing member, and the injection port is encapsulated by a resin after the oil injection. For high-speed open/close operation of the shutter, oil having low viscosity is desired, but in this case, the speed of oil when passing through the injection port increases to cause a problem of damaging the shutter."}
-{"text": "Modern communication service providers frequently offer a wide variety of communications services to users. For example, users often purchase, subscribe to, or otherwise obtain voice calling, text messaging, or high speed data access services from communication service providers. In addition, a service provider may claim a certain level of data transfer speed when marketing or branding a high speed data access service.\nA service provider making such a speed claim with regard to an offered high speed data service may desire verification of the actual experienced data rates of a user. While the service provider may monitor certain network conditions from the perspective of the communication network, such measurements merely provide a rough estimate of actual user experience. Unfortunately, existing methods to determine user experience, such as drive testing the network or providing customer satisfaction surveys, are not exhaustive enough to capture a significant data set that accurately represents the entire network."}
-{"text": "As disclosed in a patent document, Japanese Patent Laid-Open No. JP 2004-326629 A, (patent document 1) listed below, an abnormality monitoring device that monitors an abnormality of the microcomputer is known. The abnormality monitoring device includes an abnormality detector that detects an abnormality during a normal operation of the microcomputer and a standby disorder detector that detects a disorder during a standby time of the microcomputer.\nThe microcomputer outputs a watch dog pulse during the normal operation, and does not output the watch dog pulse during the standby time. Therefore, the abnormality detector determines that the microcomputer in the normal operation has an abnormality caused therein when the watch dog pulse is not input to the abnormality detector from the microcomputer for more than a preset period of time. The standby disorder detector determines that the microcomputer has a disorder caused therein when the watch dog pulse is input to the standby disorder detector from the microcomputer during the standby time of the microcomputer.\nAs described above, the standby disorder detector disclosed in the patent document 1 determines that the microcomputer has a disorder caused therein when the watch dog pulse is input to the standby disorder detector from the microcomputer during the standby time of the microcomputer. However, depending on a state of disorder of the microcomputer in the standby time, the microcomputer may possibly not output the watch dog pulse during the standby time (i.e., designated as a low power consumption mode, or a low power mode). Therefore, the abnormality monitoring device in the patent document 1 is not capable of detecting such a disorder of the microcomputer."}
-{"text": "Vascular endothelium is a cellular layer lining the inner part of blood vessels, including arteries and veins. Endothelium is presently looked upon as an important metabolically active autocrine/paracrine/endocrine organ that regulates cardiovascular function and maintains vascular homeostasis by: modulating vascular tone; regulating solute transport into cell components of the vessel wall, local cellular growth, and extracellular matrix deposition; protecting the vessel from the potentially injurious consequences of substances and cells circulating in blood; and regulating the hemostatic, inflammatory, and reparative responses to local injury. One of the main functions of endothelium is to produce or release substances, such as nitric oxide (NO), that control the behaviour of the blood vessels such as their dimensions, elasticity, permeability and reactivity, including the ability to constrict and dilate. Endothelium-derived mediators regulate not only blood flow and permeability vascular elasticity, reactivity and structure, but also local and systemic inflammatory response as well as thromboresistance of vessels. Vasoprotective endothelial mediators such as nitric oxide (NO), prostacyclin (PGI2) endothelium-derived hyperpolarising factor (EDHF), bradykinin (Bk) tissue plasminogen activator (t-PA), thrombomodulin (TM) or ADP-ase do exert antithrombotic, anti-inflammatory and vasoprotective action.\nOn the other hand, excessive production by endothelium of superoxide anions (O2-), isoprostanes, angiotensin II (ang II), endothelin 1 (ET-1), plasminogen activator inhibitor (PAI-1), tissue factor (TF), von Willebrandt factor (vWF), chemokines (e.g. monocyte chemotactive protein MCP-1), cytokines (e.g. IL-6), and increased expression of adhesion molecules (e.g. selectin P, ICAM-1) promote inflammation and thrombosis of vascular wall that may eventually lead to the development of atherosclerotic lesion. Accordingly, healthy endothelium is essential for undisturbed functioning of the cardiovascular system, while endothelial dysfunction leads to its various pathologies. In particular, endothelial dysfunction is pivotal to atherogenesis, it is present at the earliest stages (e.g. preceding angiographic or ultrasonic evidence of obstructive plaque) as well as later stages of arterial disease, contributing to clinical sequelae related to tissue damage (eg, ischemia, infarction, and organ failure).\nEndothelial dysfunction in most general terms refers to an impairment of the ability of the endothelial cell layer to produce an appropriate vasodilatory response to stimuli. Many studies provided evidence that endothelial dysfunction (assessed on the basis of the impairment of NO-dependent vasodilatation) may be regarded as prognostic factor for the development of adverse cardiovascular events. Indeed, relative risk for adverse outcomes is elevated approximately 10-fold when there is evidence of coronary or peripheral endothelial dysfunction.\nVarious conditions, including hypercholesterolemia, systemic hypertension, smoking, diabetes, congestive heart failure, pulmonary hypertension, estrogen deficiency, hyperhomocysteinemia, and the aging process itself, have been associated with impaired function (dysfunction) of endothelium. As a result, the vessel wall in these conditions may promote inflammation, oxidation of lipoproteins, smooth muscle proliferation, extracellular matrix deposition or lysis, accumulation of lipid-rich material, platelet activation, and thrombus formation. All of these consequences of endothelial dysfunction may contribute to development and clinical expression of atherosclerosis. The potential consequences of endothelial dysfunction further include coronary constriction or inadequate dilation during physical or mental stress, producing myocardial ischemia; plaque rupture and thrombosis, causing unstable angina or myocardial infarction; and reperfusion injury after thrombolysis.\nSeveral methods and apparatuses for non-invasive evaluation of the health of vascular endothelium in vivo have been developed.\nIn particular, methods are known that are based on monitoring the physiological conditions or characteristics of the arteries in the patient's limb after reactive hyperemia.\nReactive hyperemia is a physiological phenomenon that occurs in a patient after blocking (or occlusion) of a major artery. Such blocking or occlusion of artery in the limb, such as brachial artery, is typically done by inflating a blood pressure cuff slightly above systolic pressure for a period of about 5 minutes. Anoxia or severe hypoxia in the limb downstream from the occluded artery is usually a result of such blocking. Sudden release of the blocking causes endothelial cells to react by generating NO and dilating. The phenomenon of reactive hyperemia lasts up to 10 minutes before return to pre-test blood volume values. Blood flow is a characteristic of the artery, and under reactive hyperemia blood flow through an artery, vein or limb is significantly greater as compared with normal blood flow.\nCurrently the most popular method is flow mediated dilatation (FMD), a non-invasive technique based on monitoring of diameter of arteries after reactive hyperemia with a two-dimensional ultrasound and Doppler ultrasound. Its results correlate well with invasive coronary endothelial testing as well as with the presence and severity of coronary atherosclerosis. This technique is described for example in a review by S. Patel. And D. S. Celermajer, Pharmacological Reports 2006, 58, suppl. 3-7. However, this method is quite expensive, requires sophisticated equipment and highly specialized operators, is highly operator dependent and is poorly reproducible due to variability of measurements and poor resolution relative to arterial size. Hence, its use is limited and the method is not applicable on a more general basis.\nFor the purpose of assessment of vascular endothelial function changes of other physical parameters in response to reactive hyperemia have been also used, such as fingertip skin temperature (WO2005118516; N. Ahmadi et al. Int. J. Cardiovasc. Imaging (2009) 25:725-738), blood pressure in a finger (pulse wave amplitude) using plethysmography (EP1360929, EP1992282, WO00/57776, EP2110074), and peripheral arterial tone (WO2000/074551, WO2002/034105).\nNon-invasive technique for detection of endothelial dysfunction based on monitoring blood flow related changes in the level of a substance present in a limb after reactive hyperemia is disclosed in WO03/051193. The method involves blocking blood flow in the limb to stimulate endothelial function and then releasing the blood flow block to observe, measure and record said changes as a function of time, said changes being indicative of endothelial dysfunction. Said substance can be a tracer substance injected in a vein, such as a radiation emitter or a contrast agent, and the ingress of said tracer into the limb is detected and measured, for example by means of gamma ray detection. Tracer measurement in a pair of two laterally opposed limbs should be performed and the tracer presence compared between both limbs. Alternatively, a physical characteristics of the limb, such as temperature or color, or a property of a metabolic or other biochemical product circulating in the limb following the release of the blood flow block, such as O2, CO2 or reduced hemoglobin, is measured by a suitable technique. As suitable techniques there were suggested gas emissions across the skin surface within a cell placed on the skin surface, optical techniques, such as spectral analyzers or optical transmission/diffusion detectors, such as the visible-reflectance hyperspectral analysis, and EPR/NMR techniques. Either the appearance rate of a depleted substance or the disappearance (depletion) of an accumulated product can be detected. A rate of change of the measured parameter shortly after release of the occlusion or blockage is suggested as a primary factor in determining endothelial dysfunction. In the case of the use of a tracer, the rate for both the blocked limb and the contra-lateral control limb is measured.\nIt has been established that endothelial dysfunction is an early event and major risk factor for atherosclerosis and an important indicator for a medical professional, allowing for early diagnosis of the risk of cardiovascular disease.\nTesting endothelial function is therefore a highly desirable alternative for a diagnostic approach based on performing a set of various biochemical tests, especially in apparently healthy individuals, i.e. individuals not showing any signs of cardiovascular disease.\nA testing method is needed that would allow to evaluate function and detect any dysfunctions at an early stage of impairment in order to identify patients for prophylactic or therapeutic intervention to improve the dysfunction and/or for further more detailed and complicated diagnostic tests.\nThe need exists to provide a non-invasive test for evaluating endothelial function which would be reliable, easy to carry out and inexpensive, and thus applicable for tests in large patient populations, for example for screening purposes.\nThere is also a need for a simple, quick and non-expensive test that would allow to monitor and control the response of a patient to a medical treatment of cardiovascular disease."}
-{"text": "Projectiles currently used for the M203/M320 gun launcher devices include the M433 High Explosive Dual Purpose Round, M406 High Explosive Round, M583A1 Star Parachute Round, M585 White Star Cluster Round, M713 Ground Marker Round, M781 Practice Round, M651 CS Round and M576 Buckshot Round. The M203 is not a stand-alone gun. It must be attached to an M16 or M4 rifle, for instance, and fires a family of 40 mm low velocity projectile rounds. In order to load a 40 mm round into an M203, the forward sleeve of the gun tube is slid forward, and a 40 mm round is pushed into the sleeve. The sleeve is then slid backwards and locked into position, ready to fire. The barrel of the M203 is rifled, which spins up the projectile as it travels down the barrel tube, and the projectile leaves the gun tube with a high spin rate.\nThe 40 mm cartridge case is typically not reused or recycled after firing. If the cartridge case becomes stuck in the gun tube, there is a tool to help the war-fighter push it from the tube. The M320 gun launcher device fires the same 40 mm low velocity ammunition as does an M203. There are several improvements that the M320 has compared to the M203. The M320 has stand alone single shot capability and may be fired by a war-fighter without attaching to a gun. It also has the capability to be attached to and fired from an M16/M4 rifle, similarly to an M203 gun launcher device. A major improvement of the M320 is the ability of the firing tube to open sideways. By opening sideways, longer ammunition can be loaded into the M320 gun tube as compared to loading into an M203. While the approximate maximum length of the projectile is approximately five inches in the M203, projectiles several inches longer can be loaded into an M320. Both M203 and M320 have rifled barrels which induce spinning in the projectile as it travels through the gun tube.\nThere exists a great need to develop a new family of 40 mm projectiles that can have extended range through increased velocity and projectile shape to over 1000 meters; can provide capability for an add-on optional guidance-navigation and control (\u201cGN&C\u201d) subsystem; can increase lethality through enhanced precision and accurate strikes; can increase projectile length which can accommodate more fragments and an increased amount of explosives and/or energetics; and can provide a novel system with much longer rear end fins, that can allow for greater stability, in applications which will no longer need any appreciable, or no, projectile spin for flight accuracy.\nCurrently, low velocity (\u201cLV\u201d) 40 mm projectiles are length restricted by the presence of the loading system that is in the M203 gun. The 40 mm LV projectiles do not fly as well as desired, due to the inability to add fins for providing stability. Instead, 40mm LV projectiles are spin stabilized, which both do not fly as well as desired and also can not employ guidance systems due to these projectiles' extensive pitch, yaw and drag. Current guidance-navigation and control (\u201cGN&C\u201d) subsystems can't function in a spin environment. Guided ammunition would need to have no spin (or very little spin) so that the seeker or camera therein, e.g., can see the target clearly while in flight. However, the conventional M203 or M320 grenade launcher guns do have spin since they have internal rifling intended to spin up the projectiles (at approximately 1800 revolutions per second) for stabilization of flight. Clearly then, a way must be found to avoid the spin which is inherent in conventional M203 or M320 grenade launcher guns when optional guidance is desired to be added to the projectile.\nRecently developed new 40 mm projectiles of several varieties provide better performance, extended range, with reduced or no spin at all, and have the ability to have added on guidance and control systems to hit stationary targets. These novel technologies all allow for better performance, yet while reducing pitch, yaw, and/or drag. The stability of these projectiles are still somewhat limited by the length and area of their fins. These projectiles will fly well and have the ability to add guidance and control with small canards added in the front of the projectile, in a canard actuating system (\u201cCAS\u201d). Smaller (width and/or length) canards however would mean less control authority or less of an ability to make large course corrections while in flight. The projectiles' center of gravity (\u201cCG\u201d) and center of pressure (\u201cCP\u201d) needs to be a certain distance apart in order to provide a needed stability margin. In the case of fin stabilized projectiles, the CG needs to be in the front of the projectile while the CP needs to be in the rear of the projectile. Larger canards, and changing the position of the canards away from the front of the projectile, move the CP forward and begin to compromise or reduce the stability margin. By incorporating longer fins on a projectile with more surface area however, the CP is moved back and the projectile gains a better stability margin. With this greatly improved stability margin, larger canards can be positioned where needed. This improvement now allows for a projectile to have larger course corrections and hit harder targets such as moving targets. Optional guidance and control can now be added to the front or side of a projectile utilized for battle. The optional guidance and control provides increases accuracy and precision hits on targets up to and beyond 1000 meters (a goal which is sought). The obturator on the projectile prevents gases from blowing past it; it seals in the gases allowing a projectile to move with a constant rate push as it travels down the barrel of a gun launcher. The obturator must be positioned immediately behind of the cartridge case. This is because the M320 and M203 have both a section to seat the cartridge case, and then also enough room still, for the obturator before the rifling part of the gun tube is encountered. The obturator is higher than (near the end of) the projectile and therefore is found at the top of the gun tube. If the obturator were made too long in dimension, then the cartridge could not be loaded. This is since the obturator cannot be pushed into the rifling just by hand. Currently, there is an approximately one inch long of space in between the back of the obturator and the front of the (inside of the) cartridge case. This can only allow for fins that are one inch long therefore. The projectile continues on past the obturator by a distance of approximately 1/16 inch (press fit into the cartridge case) and then slopes down at an angle of 15 degrees, toward the end of the projectile. There is currently provided a housing for the fins at the end of the projectile, which housing functionally attaches the fins to the projectile. The fins physically now extend from the housing and up until the beginning of the interference fit with the cartridge case. Therefore the fins are approximately one inch long. The fins open up after exit from the gun. These finned projectiles have tested with excellent flight results, however, to add in a guidance and control (G&C) canard system as desired, the fins would not allow for larger and more aggressive canards and positioning. These projectiles are good for stationary targets. To hit moving and stationary targets, fins that are two inches or longer are needed. Therefore, what is greatly needed is to provide a low velocity 40 mm projectile with two or more inch long fins to allow for aggressive guidance and control CAS systems. In addition, aeroballistic modeling shows that the rear end of the projectile will have less drag with slopes on the projectile that are less than eight degrees, rather than the currently mentioned fifteen degrees, and this type of projectile is also greatly needed."}
-{"text": "Workstations with one or more robots are known from practice in various forms. As a rule, they are designed as fully automatic stations and they are equipped with protective partitions, which mechanically prevent the access of workers or switch off the industrial robot, for accident prevention for the workers.\nFurthermore, there are efforts to allow humans to cooperate or collaborate with industrial robots, especially tactile robots. This is called human-robot cooperation (HRC). Tactile articulated arm industrial robots suitable for this are known from, e.g., DE 10 2007 063 099 A1, DE 10 2007 014 023 A1 and DE 10 2007 028 758 B4.\nOn the other hand, a safety monitoring at an industrial robot by means of desired values is known from DE 10 2007 059 481 A1. DE 10 2005 027 522 B4 discloses an industrial robot with a moving monitored area around the tool or the dangerous area by means of a capacitive sensor system. A risk assessment according to DIN EN ISO 12100 is to be performed for designing workstations."}
-{"text": "1. Field of the Invention\nThe present invention relates to a drill bit sharpener, and more particularly to a drill bit sharpener having a rotating grinding wheel with an adjustable position to sharpen drill bits, which can be assembled quickly.\n2. Description of Related Art\nDrilling machines and drills make holes in objects. A drill has a drill bit that is turned to make the hole in the object. A cutting edge with a relief angle is formed on one end of the drill bit. After an extensive use, the cutting edge and the relief angle of the drill bit will become dull or break. A user must replace the dull or broken drill bit with a new drill bit to keep making holes efficiently, but replacement of the drill bits is not economical.\nTo sharpen a dull or broken drill bit is more economical than throwing the drill bit away. A conventional bench grinder with a rotating grinding wheel is often used to sharpen drill bits. The user holds the drill bit in his or her hands and uses the rotating grinding wheel to grind a cutting edge with the appropriate relief angle on the drill bit. Using a conventional bench grinder to sharpen drill bits has some shortcomings in both operating and using. The shortcomings include the following.\n1. Inconvenience operation:\nBench grinders are big and occupy a big space. The bench grinder must be mounted in a location before it can be used. The bench grinder is not easy to carry and not portable. Furthermore, holding the drill bit at the proper angle against the grinding wheel to accurately the tip is not easy. A user often much time to grind proper shape and angle on the drill bit.\n2. Hazardous operation:\nThe grinding wheel rotates at a very high speed, and at least part of the grinding wheel is exposed to surrounding environment. The user holds the drill bit in his or her hands and may inadvertently touch the grinding wheel if he.or she is the least bit careless. The high-speed, grinding wheel will injure the user if the user touches the wheel.\nTo overcome the shortcomings, the present invention provides a drill bit sharpener to mitigate or obviate the aforementioned problems.\nThe main objective of the invention is to provide a drill bit sharpener that will sharpen drill bits easily, quickly and efficiently.\nAnother objective of the invention is to provide a drill bit sharpener that is portable so a user can carry the sharpener and use it on the work site.\nOther objectives, advantages and novel features of the invention will become more apparent from the following detailed description when taken in conjunction with the accompanying drawings."}
-{"text": "1. Field of the Invention\nThe present invention relates to a call device, a call control system, a call management system, and a call control method.\n2. Description of the Related Art\nThere have heretofore been available in the art cellular phone terminals which output a melody (hereinafter referred to as \u201cringer melody\u201d) that is preset by the user of a cellular phone terminal as a recipient's terminal during a period of time after the cellular phone terminal has received an incoming call until the user starts to talk to the calling party. There have also been available in the art cellular phone terminals which output a melody (hereinafter referred to as \u201cring back melody\u201d) that is preset by the user of a cellular phone terminal for a recipient's terminal during a period of time after a cellular phone terminal has sent an outgoing call until the parties start talking to each other.\nThe user of a cellular phone terminal presets a preferred melody as a ringer melody or a ring back melody. Therefore, the user, who owns a cellular phone terminal which outputs a ringer melody or a ring back melody, can enjoy the music until the parties start talking to each other through the cellular phone terminal.\nThere has been proposed a call device, which is capable of superimposing a music signal on a speech signal representative of speech for the user to enjoy the music while talking to the other party.\nFor example, Document 1 (Japanese laid-open patent publication No. 2002-101206) discloses a telephone communication service management device for superimposing speech advertisement information such as music information on speech information and outputting the combined information. The telephone communication service management device presets the speech advertisement information.\nDocument 2 (Japanese laid-open patent publication No. 2002-344571) reveals a call device for superimposing a music signal on a speech signal and outputting the superimposed signal. The user of a cellular phone terminal presets the music signal.\nRinger melodies or ring back melodies give a certain level of enjoyment to the users of cellular phone terminals.\nHowever, it is desirable for ringer melodies or ring back melodies to give further enjoyment to the users of cellular phone terminals.\nAt present, a cellular phone terminal outputs a ringer melody or a ring back melody for a short period of time from the arrival of an incoming call to the start of the connection. Therefore, the user of a cellular phone terminal can enjoy the music that is preset by the user of a cellular phone terminal as a recipient's terminal, for a short period of time from the arrival of an incoming call to the start of the connection.\nAccording to Document 1, the music and the speech which are outputted together, are provided by the telephone communication service management device, without any reference to the called party. Therefore, the telephone communication service management device may possibly output music together with speech, which the called party is not fond of.\nThe telephone communication service management device disclosed in Document 1 is also disadvantageous in that if a ringer melody or a ring back melody is preset, then the music, which is output after the arrival of an incoming call until the parties start to talk to each other, is different from the music, which is output while the parties are talking to each other. Therefore, the user of a cellular phone terminal can enjoy the music (a ringer melody or a ring back melody) that is preset by the user of a cellular phone terminal for a recipient's terminal, only during a short period of time after the arrival of an incoming call until the connection is established.\nWith the call device disclosed in Document 2, if a ringer melody or a ring back melody is preset, then the music, which is output from the arrival of an incoming call until the parties start to talk to each other, may possibly be different from the music, which is output while the parties are talking to each other. Consequently, the user of a cellular phone terminal can enjoy a ringer melody or a ring back melody, only during a short period of time after the arrival of an incoming call until the connection is established.\nAnother problem is that the called party who has preset a ring back melody is unable to confirm the amount of the ring back melody to which the calling party has already listened. Therefore, the called party cannot judge whether the calling party has fully enjoyed the ring back melody or not."}
-{"text": "The present invention relates to an assembly of a buckle mechanism and blade, and a method of securing the blade to and releasing the blade from the buckle mechanism. More specifically, the present invention relates to a buckle mechanism having a locking surface, a movable guide, and a unitary biasing member.\nExisting buckle mechanisms have a moving lock member that secures a blade. These known buckle mechanisms are very complex devices comprised of a plethora of parts designed to engage and release the moving lock member from the blade to accomplish its relatively simple purpose. As a result of their complexity, known buckle mechanisms with a moving lock mechanism can be relatively expensive to manufacture and may be susceptible to failure.\nOne of the failures associated with these known buckle mechanisms is that releasing the blade from the buckle mechanism is difficult, if not impossible, after the assembly has withstood particular tensile forces. For example, vehicle passengers, after having been involved in an accident and restrained by a seat belt system employing a buckle mechanism with a moving lock, are unable to release the buckle mechanism. It is believed that the tensile forces applied to the moving lock member prevent the lock member from being moved in these conditions. Thus, a need exists for an assembly of a buckle mechanism and blade that has a minimum number of parts and avoids the difficulties of a moving lock member in a buckle mechanism.\nAccordingly, the present invention is directed to an assembly of a buckle mechanism and blade, having a fastening surface, including a base, having a surface and a projection extending away from the surface, and a guide supported by the base. The projection has a lock surface substantially perpendicular to a plane passing through the surface of the base and contiguous to the surface. The guide engages the fastening surface of the blade with the lock surface of the projection.\nIn an alternative embodiment, the buckle mechanism and blade assembly includes a base within the housing, a guide supported by the base, and at least one biasing member within the housing. The biasing member locates the guide to a first position in which the blade engages the projection and opposes movement of the guide to a second position in which the blade disengages from the projection. The base has a surface and a projection extending away from the surface, and the guide locates the blade proximate the surface.\nIn a further embodiment, the invention includes a buckle mechanism having a base with a projection and a guide supported by the base. The guide has at least one wall defining a void and is positionable in a first position so that the projection occupies at least a portion of the void and a second position so that the projection is removed from the void.\nThe present invention also provides a method of securing a blade to a buckle mechanism. In this method, the buckle mechanism has a base and a guide supported by the base, the base has a projection providing a lock surface and a portion aligned with the lock surface. The guide is configured to support the blade, and the blade has a fastening surface and a support surface adjacent the fastening surface. This method includes: (1) inserting the blade into the guide; and (2) engaging the fastening surface of the blade with the lock surface of the projection and the support surface of the blade with the portion of the base aligned with the lock surface.\nAn alternative method of the present invention is a method of releasing a blade from a buckle mechanism, where the buckle mechanism has a base and a guide supported by the base. The base has a surface and a projection extending away from the surface. The projection provides a lock surface, and the guide is configured to support the blade. The blade has a fastening surface and a support surface adjacent the fastening surface. This method includes: (1) moving the fastening surface of the blade by the guide in a first direction substantially oblique to a plane passing through the surface of the base to disengage the fastening surface of the blade from the lock surface of the projection; and (2) removing the blade from the guide in a second direction substantially parallel to the surface of the base.\nA further method of the present invention is a method of securing a blade to a buckle mechanism and releasing a blade from the buckle mechanism. The buckle mechanism has a base and a guide supported by the base. The base has a surface and a projection extending away from the surface, where the projection provides a lock surface. The guide is configured to support the blade, which has a fastening surface. This method includes: (1) inserting the blade into the guide; (2) engaging the fastening surface of the blade with the lock surface of the projection; (3) moving the fastening surface of the blade by the guide in a first direction substantially oblique to a plane passing through the surface of the base to disengage the fastening surface of the blade from the lock surface of the projection; and (4) removing the blade from the guide in a second direction substantially parallel to the surface of the base."}
-{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing a device for performing a dicing in such a way that a substrate, on which two or more devices and alignment marks for positioning are formed, is positioned in accordance with the alignment mark, a positioning method, and dicing method and dicing apparatus in which the substrate is separated through the dicing to individual device elements equipped with individual devices.\n2. Description of the Related Art\nHitherto, to perform a dicing for a device substrate in which two or more devices, such as semiconductor wafers, are made, generally, there is adopted such a process that the device substrate is fixed on a dicing table through putting the device substrate on an adhesive dicing tape, the device substrate is subjected to the dicing processing so as to be segmented into device elements referred to as a chip, and the device elements are picked up one by one from the dicing tape.\nHowever, a state-of-the-art device including MEMS (Micro Electro Mechanical Systems) device, on which the research is advanced, has an extremely minute structure. Thus, when the device elements are picked up from the adhesive tape, the structure is damaged. This involves a problem that the yield is extremely low.\nIn order to improve this problem, there is proposed a method in which a device substrate is wrapped with the coagulant, without the use of the adhesive tape, and the coagulant is coagulated, so that the device substrate is fixed (cf. Japanese Patent Publication TokuKai. Hei. 10-230429, and Japanese Patent Publication TokuKai. Hei. 11-309639).\nHowever, fixing of the device substrate by the coagulation of the coagulant involves a problem of a difficulty of positioning of the device substrate. For the dicing for the device substrate, there is a need that the device substrate is exactly positioned relatively to a dicing blade. To perform the positioning, when devices are made on the device substrate, an alignment mark is formed at a position determined exactly as to the devices, and the alignment mark is photographed by a camera, so that the positioning is carried out based on the alignment mark on the image. However, when the coagulant is coagulated, the coagulation brings about a state that the coagulant is opaque or light scattering. Thus, it is difficult for a camera to recognize the alignment mark. Before the coagulation of the coagulant, it is easy for the camera to recognize the alignment mark. However, the position may be shifted in the process in which the coagulant is coagulated, and thus the recognition of the alignment mark before the coagulation would not be useful for positioning.\nIt is considered that the device substrate is fixed with a material liquid with viscous high in place of the coagulant. In this case, it is possible to recognize the alignment mark. However, as compared with a case where a position is fixed by coagulation of the coagulant, the position of the device substrate changes easily. And, the dicing for the device substrate is carried out while the coolant such as the pure water is sprayed for cooling and the removal of the cutting powder. However, even if the device substrate is covered with a liquid material of high viscosity, the device substrate is pressured by spraying of the coolant, and as a result, the device will be damaged. This involves a problem that the yield is not sufficiently improved.\nJapanese Patent Publication TokuKai. Hei. 6-331813 discloses a technology in which the substrate that forms the diffraction grating is put in the water tank that fills the pure water, and the pure water is frozen to fix the substrate and then subjected to cutting. This also easily brings about opaque or light scattering state at the time of freezing. Thus, according to the technology disclosed in Japanese Patent Publication TokuKai. Hei. 6-331813 too, it is difficult to perform a positioning base on the alignment mark.\nFurther, according to the second type of extended library apparatus show in FIG. 21, there are provided the same number of moving mechanisms as the number of coupled library apparatuses. Thus, as compared with the first type of extended library apparatus show in FIG. 19, which is capable of coping with coupling of a plurality of library apparatuses with one moving mechanism, it involves problems of high cost and increment of dissipation power."}
-{"text": "With a conventional radiographic apparatus, an X-ray source projects X-rays to an object (e.g., a medical patient). X-ray beams transmitted through the object are detected to radiograph the object by a screen film cassette, film autochanger, CR (Computed Radiography), FPD (Flat Panel Detector), or the like.\nIn the field of radiography, a high-resolution solid X-ray detector which uses an FPD is proposed. This solid X-ray detector has an X-ray sensor including a two-dimensional array in which arrays of 3,000 to 4,000 photoelectric conversion devices (e.g., photodiodes) are arrayed two-dimensionally. Each photoelectric conversion device generates an electrical signal corresponding to the X-ray dose projected to the X-ray sensor. The X-ray image of an object is obtained by arranging the object between an X-ray source and X-ray sensor and converting the X-ray dose which has been transmitted through the object into an electrical signal. A signal from each photoelectric conversion device is read out individually and digitized, and thereafter image-processed, stored, or displayed.\nFIG. 9 is a conceptual view showing the structure of a system which includes a conventional cassette type radiographic apparatus. As shown in FIG. 9, a radiographic apparatus 801 incorporates an X-ray detector 802. X-rays generated by an X-ray generator 803 irradiate an object 804. The X-rays transmitted through the object 804 are detected by photoelectric conversion devices (not shown) which are arrayed like a matrix on the X-ray detector 802. Electrical signals output from the photoelectric conversion devices are image-processed by an image processor 805 to display the X-ray image of the object 804 on a display 806 such as a monitor.\nIn recent years, a low-profile, higher-density mounting technique has improved, and a compact, low-profile solid X-ray detector which uses an FPD is becoming possible, e.g., an X-ray screen film cassette (see Japanese Patent Laid-Open Nos. 2003-057352 and 2002-186604).\nFIG. 10 is a side sectional view showing an example of a cassette type radiographic apparatus which uses an FPD. As shown in FIG. 10, an electronic cassette used for radiography or the like has a particle phosphor 131, e.g., GOS, which converts X-rays into visible light, a MIS photosensor portion 115 which uses amorphous silicon and photoelectric conversion devices 109 which are arranged like a matrix to convert the visible light into electrical signals, a TFT switching portion 116, a base 110, a circuit board 111 which supports the base 110, a circuit board 113 on which electronic components for processing the photoelectrically converted electrical signals are mounted, wiring lines 114, case lid 101 which is used to house the above members, and a case main body 117. A buffer member 102 serving as a relaxation portion that relaxes a force from outside the housing is arranged between the case lid 101 and particle phosphor 131. A resin 130 made of PET or the like is arranged between the buffer member 102 and particle phosphor 131.\nThe particle phosphor 131 is adhered to the photoelectric conversion devices 109 by an adhesion layer 106 through a second protection layer 107 made of an organic substance such as PI and a first protection layer 108 made of a nitride or the like. The circuit board 111 on which the electronic components 113 for processing the electrical signals photoelectrically converted by the photoelectric conversion devices 109 are mounted is planarly attached to the lower surface of the base 110 in tight contact through an insulating sheet 112.\nConventionally, in the electronic cassette which uses the particle phosphor 131 such as GOS, the photoelectric conversion devices 109 which are made of a component such as glass are more vulnerable to an external force than the particle phosphor 131 and protection layers 107 and 108. Therefore, a cassette type radiographic apparatus which uses an FPD has been designed with reference to the strength of the photoelectric conversion devices 109.\nIn the conventional cassette type radiographic apparatus which uses the FPD, problems occur when a columnar crystal phosphor such as CsI (cesium iodide crystal) is used in place of the particle phosphor 131. The reasons of the problems are roughly classified into two. According to the first reason, the columnar crystal phosphor is fractured by a weaker external force than the photoelectric conversion devices made of a component such as glass. As the columnar crystal phosphor is more vulnerable to the external force, the conventional technique for protecting the photoelectrical conversion devices cannot sufficiently protect the columnar crystal phosphor. According to the second reason, as the columnar crystal phosphor has stricter demands for an external force acting on a small area than a granular phosphor used in the conventional cassette type radiographic apparatus using the FPD, the crystals of the columnar phosphor may be broken. When a stress is applied, holes may be formed in the protection film by steps formed by variations in crystal length.\nAccording to the prior art, a standing- or lying-position radiographic apparatus is available which uses a columnar crystal phosphor and an FPD. When this apparatus is used as a portable cassette type radiographic apparatus, problems occur. The reasons of the problems are roughly classified into two. According to the first reason, unlike the standing- or lying-position portable cassette type radiographic apparatus, the portable cassette type radiographic apparatus is used in various applications, and accordingly the weight of the object may be applied to it. Therefore, a stress absorbing portion which is not necessary in the standing- or lying-position radiographic apparatus is necessary in the portable cassette type radiographic apparatus in case various types of external forces are applied when, e.g., the operator erroneously hits the apparatus with something from above or places his or her elbow on the apparatus. According to the second reason, the cassette type radiographic apparatus is not provided with a grid space. The standing- or lying-position radiographic apparatus usually uses a scattered ray removing mechanism (grid), and the scattered ray removing mechanism (grid) serves as a stress absorbing portion. In the cassette type radiographic apparatus, however, the radiographic target is often a body portion, e.g., a limb, which does produce many scattered rays, and accordingly a scattered ray removing mechanism (grid) is not usually mounted due to the requirements for a lower profile and lighter weight. Therefore, a structure is necessary which protects the columnar crystal phosphor from the external force.\nIn the cassette type radiographic apparatus, for achieving a lower profile., a gap must be minimized as small as possible in the direction of thickness of the cassette. With a small gap, however, if something is erroneously dropped on the case to deform it elastically, the phosphor may be broken at a high possibility. In particular, in the case of a columnar crystal phosphor such as CsI, when the stress acts on the crystals to break the phosphor crystals, the directional characteristics of light in the phosphor change, and a scar may be undesirably present in an image obtained by X-ray irradiation.\nFIGS. 11A to 11D show images obtained when a stress acts on the columnar crystal phosphor to break the phosphor crystals. After an experiment (FIG. 11A) of dropping a screwdriver onto the columnar crystal phosphor (CsI: TI+) was conducted, the columnar crystal phosphor was adhered to photoelectric conversion devices, and radiography was performed (FIG. 11B). As shown in FIG. 11B, a large hitting mark ranging for a diameter of 3 mm appears on the image. FIG. 11C is a graph showing the sectional profile of an output value of a photodetector taken along the line A\u2013A\u2032 of FIG. 11B. As shown in FIG. 11C, a portion where the optical output increases by about 10% and a portion where the optical output decreases by about 10% are present. FIG. 11D Is a graph showing the sectional profile of the output value of the photodetector taken along the line B\u2013B\u2032 of FIG. 11B. As shown in FIG. 11D, many portions are obviously present where the optical output increases by about 10%. In this manner, with reference to FIGS. 11C and 11D, portions where the output value Increases or decreases are obviously present, and a difference in output value is as large as \u00b110%. The output value fluctuates probably because the phosphor crystals are broken and a light-outputting portion changes. As is apparent from the experimental images of. FIGS. 11A to 11D, since the columnar crystal phosphor has a low strength, even when only a lightweight material such as a screwdriver is erroneously dropped, it forms a scar In the image to cause a clinical problem."}
-{"text": "The difficulties in the development of hair coloring compositions, which can deliver precise long-lasting colors with low hair damage, are in part due to the inherent structure of the hair itself and in part due to the necessary harsh conditions of effective hair coloration processes.\nIn general, the condition and structure of human hair is not regular along the length of the hair shaft. Human hair is subject to various chemical and mechanical treatments such as combing, brushing, shampooing, heating, perming as well as exposure to the sun. As such, the hair at the ends of the hair shaft will generally exhibit signs of damage relative to the new growth close to the scalp. This damage can lead to inconsistent coloration when the hair is dyed due to irregular uptake of the hair-coloring agents along the length of the hair shaft. Thus there is a need for hair coloring compositions which can deliver substantially consistent hair color results throughout the hair.\nOnce the hair has been colored there is a desire for the color to be resistant to the fading action of the sun and other exterior factors and that the color is retained in a consistent manner for a predictable period of time. Thus it would be desirable to develop a hair coloring composition which exhibited reduced fade and provided improved resistance to wash out during regular cleansing regimen. An additional difficulty commonly associated with the dyeing of human hair is the need for dye systems that avoid any adverse effect on the hair and skin of the user.\nTo color human or animal hair using conventional oxidative dye technology it is generally necessary to treat the hair with a mixture of suitable oxidative coloring agents and at least one oxidizing agent and swelling agent at alkaline pHs. Hydrogen peroxide is the most commonly used oxidizing agent. This can lead to undesirable hair damage. These undesirable effects are in part due to the necessary conditions of conventional peroxide treatment, as part of the hair coloring process, which requires high pH ( greater than pH 9), use of swelling agents such as ammonia, extended exposure (from 30 to 60 minutes) and relatively high concentration of oxidizing solutions (up to 30% volume of oxygen) in order to deliver effective hair dyeing and hair melanin bleaching. Thus there is a need for hair coloring compositions, which can oxidize dye precursors and color the hair effectively without damage.\nOxidative dyes and peroxide oxidizing agents can be used to deliver a variety of hair colors to the hair. However such conventional hair coloring compositions comprising oxidative dyes and peroxide oxidizing agent do not deliver against the key consumer needs of color saturation, color vividness, precise initial color predictability, improved wash fastness, improved safety and low hair damage.\nApplicants have now found that the combination of peroxide oxidizing agents with certain transition metal salts or complexes and one or more oxidative coloring agents in hair coloring compositions can deliver excellent initial hair color in a faster time, good wash fastness of the hair color over time, desirable color saturation and vividness attributes, and significantly reduced hair damage. Furthermore, applicants have found that the combination of these dye precursors and transition metal salts or complexes with a source of peroxide at particular levels and ratios in the coloring compositions of the present invention can deliver these excellent hair coloring attributes results with minimal hair damage in a faster time.\nThus it is an object of the present invention to provide improved hair coloring compositions.\nThe following publications relate to the field of the invention:\nU.S. Pat. No. 6,004,355 discloses hair coloring compositions which comprise (a) a water-soluble peroxygen oxidizing agent, (b) an organic peroxyacid precursor oxidizing aid, and (c) one or more oxidative hair coloring agents. The compositions may be used in processes for coloring human or animal hair and may be conveniently provided in kits wherein each of the aforementioned components are individually packaged. This patent, at column 23, lines 20-35, discloses as a ligand, a list which includes 1,4,7-trimethyl-1,4,7-triazacyclononane.\nWO 97/24108, WO 97/24107 and U.S. Pat. No. 6,004,355 disclose a hair bleaching composition comprising: a) a water-soluble peroxygen bleach and (b) a bleaching aid selected from organic peroxyacid bleach precursors and/or preformed organic peroxyacids. The products can provide hair bleaching and in-use efficacy benefits including reduced hair damage at lower pH.\nU.S. Pat. No. 5,873,910 discloses a process for the two-stage oxidation dyeing of keratin fibers by applying to the keratin fibers:\nin a first stage, at least one composition A containing at least one manganese salt and/or a manganese complex,\nin a second stage, at least one composition B having a pH of greater than or equal to 6, and resulting from the extemporaneous mixing of at least one composition B1 containing at least one 4-substituted 1-naphthol and at least one composition B2 containing at least one oxidizing agent, and corresponding multi-compartment dyeing kit.\nThis patent discloses 1,4,7-trimethyl-1,4,7-triazacyclononane as a manganese complex to be used in its process at column 10, lines 20-25.\nDE 19852972 discloses the use of 1,4,7-trimethyl-1,4,7-triazacyclononane in hair dye formation.\nThe present invention is related to a hair coloring composition suitable for the treatment of human or animal hair. According to one aspect of the present invention, there is provided a hair coloring composition comprising a first composition which comprises:\n(a) a dye forming transition metal complex; which is first applied to the hair;\nand a second composition which comprises the following two compositions which are mixed just prior to application to the hair:\n(a) a water-soluble peroxygen oxidizing agent; and\n(b) one or more oxidative hair coloring agents selected from the group consisting of an aromatic diamine, an amino phenol, a polyhydric phenol, a catechol and mixtures thereof.\nAs used herein, % means weight % of the total composition unless otherwise designated. The compositions of the invention are made using known ingredients or with ingredients analogous to those known in the art. The packages or containers to be used with the compositions of the invention are made using known processes and materials or by processes and materials which are analogous to those known in the art.\nCompositions of the invention are used in a two step process for the coloring and bleaching of hair. That is, the first recited composition comprising dye forming transition metal catalysts is contacted with the hair, and then the second composition is made by mixing the following two ingredients just prior to application to the hair; (a) a water-soluble peroxygen oxidizing agent; and (b) one or more oxidative hari coloring agents. After mixture, this composition is then applied to hair. The coloring reaction takes place and the hair is rinsed.\nAs used herein the term xe2x80x9chairxe2x80x9d to be treated may be xe2x80x9clivingxe2x80x9d i.e. on a living body or may be xe2x80x9cnon-livingxe2x80x9d i.e. in a wig, hairpiece or other aggregation of non-living fibres. Mammalian, preferably human hair is preferred. However wool, fur and other melanin containing fibres are suitable substrates for the compositions according to the present invention.\nThe hair coloring compositions can contain, in addition to a mixture of active oxidizing agents, transition metal salts or complexes and oxidative coloring agents such as, by way of example, sequestrants, thickeners, buffers, carriers, surfactants, solvents, antioxidants, polymers non-oxidative dyes and conditioners.\nAs noted above, there is provided a hair coloring composition comprising a first composition which comprises:\n(a) a dye forming transition metal complex or salt;\nwhich is first applied to the hair; and a second composition, which comprises the following two compositions which are mixed just prior to application to the hair:\n(a) a water-soluble peroxygen oxidizing agent; and\n(b) one or more oxidative hair coloring agents selected from the group consisting of an aromatic diamine, an amino phenol, a polyhydric phenol, a catechol and mixtures thereof.\nThe invention also relates to a composition comprising a transition metal wherein the transition metal salt and/or transition metal complex is selected from group 1B, 6B, 7B or 8B metals and an azamacrocycle;\na composition wherein a water-soluble peroxygen oxidizing agent is selected from hydrogen peroxide, sodium perborate, sodium percarbonate, and urea peroxide complexes;\na composition wherein the transition metal complex is formed from a manganese salt and 1,4,7-trimethyl-1,4,7-triazacyclononane;\na composition wherein said manganese complex is bis(octahydro-1,4,7-trimethyl-1H-1,4,7-triazinone-N1, N4, N7)-tri-mxcexc-oxodimanganese bis hexaflurophosphate;\na composition wherein the weight ratio of water-soluble peroxygen oxidizing agent to transition metal complex is in the range of from about 6000:1 to about 1:2;\na composition wherein the weight ratio of water-soluble peroxygen oxidizing agent to transition metal complex is from about 2000:1 to about 1:1;\na composition wherein the weight ratio of water-soluble peroxygen oxidizing agent to transition metal complex is from about from 600:1 to about 60:1;\na composition wherein the said transition metal salt and/or transition metal catalyst are present in a concentration ranging from 0.0001% to 0.2% by weight of metal equivalents relative to the total weight of composition;\na composition wherein the said transition metal salt and/or transition metal catalyst are present in a concentration ranging from 0.001% to 0.1% by weight of metal equivalents relative to the weight of the total composition;\na composition which further comprises a buffering agent;\na composition having one or more oxidative hair coloring agents as described above which has a pH of about 5 to about 11;\na composition as described above which further comprises a surfactant selected from the group consisting of anionic, nonionic, cationic, zwitterionic, amphoteric surfactants and mixtures thereof; and\na hair coloring kit comprising an individually packaged oxidizing component with additional agents, an individually packaged component having one or more hair coloring agents and an individually packaged transition metal salt and/or transition metal complex with additional agents which when mixed forms a composition of the invention.\nAccording to a further aspect of the present invention, there is provided a hair coloring composition comprising:\n(b) from about 0.0001% to about 2% by weight of a transition metal and a second composition which comprises the following two compositions, which are mixed just prior to application to the hair:\n(a) a composition comprising from about 0.01% to about 10% by weight of a water-soluble peroxygen oxidizing agent;\n(b) a composition comprising from about 0.0001% to about 7% by weight of one or more oxidative hair coloring agents;\nThe invention also relates to compositions for coloring human or animal hair wherein the hair coloring composition comprises:\n(a) a composition comprising from about 0.001% to about 1.0% by weight of a transition metal salt and/or transition metal complex;\nwhich is first applied to the hair; and a second composition, which comprises the following two compositions which are mixed just prior to application to the hair:\n(a) a composition comprising from about 0.5% to about 20% by weight of a water-soluble peroxygen bleaching compound\n(b) a composition comprising from about 5% to 50% additional agents selected from surfactants, buffering agents, oxidation bases and/or couplers from about 20% to about 95.5% by weight of an inert diluent\nThe invention also relates to a process for coloring human or animal hair wherein the hair coloring composition comprises applying to the hair the compositions described above.\nIt has been found that compositions of the present invention are unexpectedly effective for lightening and coloring very dark brown or black hair.\nFor example, when (a) a dye forming transition metal complex composition of the present invention is first applied to very dark brown hair or black hair;\nand then a second composition which comprises (a) a water-soluble peroxygen oxidizing agent; and (b) a composition comprising one or more oxidative hair coloring agents the following two compositions (which are mixed just prior to application to the hair) are applied to said very dark hair brown or black hair there is obtained a very dramatic change color of the hair.\nChange in hair color can be measured by means which are known in the art. A customary scale for measuring the change in hair color is defined in terms of the variables: L, a, and b wherein L, a, and b are defined as follows: L indicates the lightness or darkness of the color value. The higher the L, therefore, the lighter the hair, and the more fading that has occurred. When L is 0, the hair is black, and when L is 100, the hair is white. xe2x88x92a and +a represent changes in color tone from green to red. xe2x88x92b and +b represent the changes in color tone from blue to yellow.\nWhat follows now is a description of the ingredients which may be used in the compositions of the invention.\nDye Forming Transition Metal Complex or Salt\nThe compositions of the invention comprise a transition metal containing catalyst for the peroxide oxidizing agents. One suitable type of catalyst is a catalyst system comprising a heavy metal cation of defined catalytic activity, such as copper or manganese cations, an auxiliary metal cation having little or no catalytic activity, such as zinc or aluminum cations, and a sequestrant having defined stability constants for the catalytic and auxiliary metal cations, particularly ethylenediaminetetraacetic acid, ethylenediaminetetra (methylenephosphonic acid) and water-soluble salts thereof. Such catalysts are disclosed in U.S. Pat. No. 4,430,243 which is hereby incorporated by reference.\nOther types of suitable catalysts include the manganese-based complexes disclosed in U.S. Pat. Nos. 5,246,621 and 5,244,594 which is hereby incorporated by reference.\nPreferred examples of these catalysts include MnIV2(xcexcxe2x88x92O)3(1,4,7-trimethyl-1,4,7-triazacyclononane) 2-(PF6)2, MnIII2(xcexcxe2x88x92O)1(xcexcxe2x88x92OAc) 2(1,4,7-trimethyl -1,4,7-triazacyclononane)2-(ClO4)2, MnIIIMnIV4(xcexcxe2x88x92O)1(xcexcxe2x88x92OAc)2-(1,4,7-trimethyl-1,4,7-triazacyclononane) 2-(ClO4)3, and mixtures thereof. Others are described in EP-A-0,549,272 which is hereby incorporated by reference. Other ligands suitable for use herein include 1,5,9-trimethyl-1,5,9-triazacyclododecane, 2-methyl-1,4,7-triazacyclononane, 2-methyl-1,4,7-triazacyclononane, 1,2,4,7-tetramethyl-1,4,7-triazacyclononane, and mixtures thereof.\nFor examples of suitable catalysts see U.S. Pat. Nos. 4,246,612 and 5,227,084 which is hereby incorporated by reference. See also U.S. Pat. No. 5,194,416 which is hereby incorporated by reference which teaches mononuclear manganese (IV) complexes such as Mn(1,4,7-trimethyl-1,4,7-triazacyclononane)(OCH3)3-(PF6). Still another type of suitable catalyst, is disclosed in U.S. Pat. No. 5,114,606, which is hereby incorporated by reference. It is a water-soluble complex of manganese (III) and/or (IV) with a ligand which is a noncarboxylate polyhydroxy compound having at least three consecutive Cxe2x80x94OH groups. Other examples include binuclear Mn complexed with tetra-N-dentate and bi-N-dentate ligands, including N4MnIII(xcexcxe2x88x92O)2MnIVN4)+ and [Bipy2MnIII(xcexcxe2x88x92O)2MnIVbipy2]-(ClO4)3.\nFurther suitable catalysts are described, for example, in EP-A0,408,131 (cobalt complex catalysts), EP-A-0,384,503, and EP-A-0,306,089 which is hereby incorporated by reference (metallo-porphyrin catalysts), and U.S. Pat. No. 4,728,455 all of which are hereby incorporated by reference.\nAlso suitable is (manganese multidentate ligand catalyst), U.S. Pat. No. 4,711,748 and EP-A 0.224,952, (absorbed manganese on aluminosilicate catalyst), U.S. Pat. No. 4,601,845; (aluminosilicate support with manganese and zinc or magnesium salt), U.S. Pat. No. 4,626,373 (manganese ligand catalyst), U.S. Pat. No. 4,119,557 (cobalt chelant catalyst) CA-A-866,191 (transition metal containing salts), U.S. Pat. No. 4,430,243 (chelants with manganese cations and noncatalytic metal cations); and U.S. Pat. No. 4,728,455 (manganese gluconate catalysts); all of which are hereby incorporated by reference.\nOther types of suitable catalysts include the manganese-based salts as described in U.S. Pat. No. 5,873,610 which is hereby incorporated by reference.\nThe compositions of the invention comprise one or more transition metal catalysts. When compared against a combination of all of the compositions of the invention after they have been mixed together (that is, the dye forming transition metal complex; the water-soluble peroxygen oxidizing agent composition; and the composition containing one or more oxidative hair coloring agents selected from the group) the concentration of each catalyst is from about 0.0001% to about 1% by weight and is preferably from about 0.001% to about 0.1% by weight.\nThe total combined level of catalysts in the compositions according to the present invention is from about 0.0001% to about 2%, preferably from about 0.001% to about 0.2%, more preferably from about 0.001% to about 0.1% by weight.\nWater-Soluble Peroxygen Oxidizing Agents\nThe compositions of the invention may comprise at least one water-soluble peroxygen oxidizing agent. Water-soluble as defined herein means a peroxygen oxidizing agent compound, which can be substantially solubilized in water.\nThe peroxygen oxidizing agents useful herein are generally inorganic peroxygen materials capable of yielding hydrogen peroxide in an aqueous solution. Water-soluble peroxygen oxidizing compounds are well known in the art and include hydrogen peroxide, inorganic alkali metal peroxides such as sodium periodate and sodium peroxide and organic peroxides such as urea peroxide, melamine peroxide, and inorganic perhydrate salt oxidizing compounds, such as the alkali metal salts of perborates, percarbonates, perphosphates, persilicates, persulphates and the like. These inorganic perhydrate salts may be incorporated as monohydrates, tetrahydrates, and the like. Mixtures of two or more such oxidizing agents can be used if desired. Preferred for use in the compositions according to the present invention is hydrogen peroxide.\nWhen compared against a combination of all of the compositions of the invention after they have been mixed together (that is, the dye forming transition metal complex; the water-soluble peroxygen oxidizing agent composition; and the composition containing one or more oxidative hair coloring agents), the peroxygen oxidizing agent is present in the compositions of the present invention at a level of from about 0.01% to about 9%, preferably from about 0.1% to about 6%, more preferably from about 0.2% to about 4% by weight.\nOxidative Hair Coloring Agents\nThe compositions of the present invention may include one or more oxidative hair coloring agents. Such oxidative hair coloring agents are used in combination with the oxidizing systems of the present invention to formulate permanent, hair dye compositions.\nPermanent hair dye compositions as described herein are compositions, which once applied to the hair are substantially resistant to wash-out. Wash-out as described herein is the process by which hair color is removed from the hair over time during normal hair cleansing regimen.\nWhen compared against a combination of all of the compositions of the invention after they have been mixed together (that is, the dye forming transition metal complex; the water-soluble peroxygen oxidizing agent composition; and the composition containing one or more oxidative hair coloring agents), the concentration of each oxidative hair coloring agent is from about 0.0001% to about 7% by weight and is preferably from about 0.001% to about 2.0% by weight.\nWhen compared against a combination of all of the compositions of the invention after they have been mixed together (that is, the dye forming transition metal complex; the water-soluble peroxygen oxidizing agent composition; and the composition containing one or more oxidative hair coloring agents), the total combined level of oxidative hair coloring agents in the compositions according to the present invention is from about 0.01% to about 15%, preferably from about 0.01% to about 10%, more preferably from about 0.1% to about 5% by weight.\nOxidative hair coloring agents which can be used in compositions of the invention can be selected from the group consisting of an aromatic diamine, an aminophenol, a polyhydric phenol, a catechol and mixtures thereof. Oxidative hair coloring agents which can also be called oxidative dyes are described in more detail below.\nThe dye forming intermediates used in oxidative dyes can be aromatic diamines, aminophenols and their derivatives. These dye-forming intermediates can be classified as; primary and secondary intermediates. Primary intermediates are chemical compounds, which by themselves will form a dye upon oxidation. The secondary intermediates, are also known as color modifiers or couplers and are used with other intermediates for specific color effects or to stabilize the color.\nThe oxidation dye intermediates which are suitable for use in the compositions and processes herein include aromatic diamines, polyhydric phenols, aminophenols and derivatives of these aromatic compounds (e.g., N-substituted derivatives of the amines, and ethers of the phenols). Primary oxidation dye intermediates are generally colorless molecules prior to oxidation. The oxidation dye color is generated when the primary intermediate is activated and subsequently joined with a secondary intermediate (coupling agent), which is also generally colorless, to form a colored, conjugated molecule.\nIn general terms, oxidation hair dye precursors or intermediates include those monomeric materials which, on oxidation, form oligomers or polymers having extended conjugated systems of electrons in their molecular structure. Because of the new electronic structure, the resultant oligomers and polymers exhibit a shift in their electronic spectra to the visible range and appear colored. For example, oxidation dye precursors capable of forming colored polymers include materials such as aniline, which has a single functional group and which, on oxidation, forms a series of conjugated imines and quinoid dimers, trimers, and the like, ranging in color from green to black. Compounds such as p-phenylenediamine, which has two functional groups, are capable of oxidative polymerization to yield higher molecular weight colored materials having extended conjugated electron systems. Color modifiers (couplers), such as those detailed hereinafter, are preferably used in conjunction with the oxidation dye precursors herein. A representative list of oxidation dye precursors suitable for use is found in Sagarin, xe2x80x9cCosmetic Science and Technology Interscience, Special Edn. Vol 2 pages 308 to 310, which is hereby incorporated by reference.\nThe typical aromatic diamines, polyhydric phenols, aminophenols, and derivatives thereof, described above as primary dye precursors can also have additional substituents on the aromatic ring, e.g. halogen, aldehyde, carboxylic acid, nitro, sulfonic acid and substituted and unsubstituted hydrocarbon groups, as well as additional substituents on the amino nitrogen and on the phenolic oxygen, e.g. substituted and unsubstituted alkyl and aryl groups.\nNonlimiting examples of suitable aromatic diamines, aminophenols, polyhydric phenols and derivatives thereof, respectively, are the following compounds:\no-phenylenediamine,\nm-phenylenediamine,\n2-nitro-p-phenylenediamine,\n1,3,5-triaminobenzene,\n2-hydroxy-p-phenylenediamine,\n2,4-diaminobenzoic acid,\nsodium 2,4-diaminobenzoate,\ncalcium di-(2,4-aminobenzoate),\nammonium 2,4-diaminobenzoate,\ntrimethylammonium 2,4-diaminobenzoate,\ntri-(2-hydroxyethyl)ammonium 2,4-diaminobenzoate,\n2,4-diaminobenzaldehyde carbonate,\n2,4-diaminobenzenesulfonic acid,\npotassium 2,4-diaminobenzenesulfonate,\nN,N-diisopropyl-p-phenylenediamine bicarbonate,\nN,N-dimethyl-p-phenylenediamine,\nN-ethyl-Nxe2x80x2-(2-propenyl)-p-phenylenediamine,\nN-phenyl-p-phenylenediamine,\nN-phenyl-N-benzyl-p-phenylenediamine,\nN-ethyl-Nxe2x80x2-(3-ethylphenyl)-p-phenylenediamine,\n2,4-toluenediamine,\n2-ethyl-p-phenylenediamine,\n2-(2-bromoethyl)-p-phenylenediamine,\n2-phenyl-p-phenylenediamine laurate,\n4-(2,5-diaminophenyl)benzaldehyde,\n2-benzyl-p-phenylenediamine acetate,\n2-(4-nitrobenzyl)-p-phenylenediamine,\n2-(4-methylphenyl)-p-phenylenediamine,\n2-(2,5-diaminophenyl)-5-methylbenzoic acid,\n2-methoxy-p-phenylenediamine,\n2,3-dimethyl-p-phenylenediamine,\n2,5-dimethyl-p-phenylenediamine,\n2-methyl-5-methoxy-p-phenylenediamine,\n2,6-methyl-5-methoxy-p-phenylenediamine,\n3-methyl-4-amino-N,N-diethylaniline,\nN,N-bis-(2-hydroxyethyl)-p-phenylenediamine,\n3-methyl-4-amino-N,N-bis-(2-hydroxyethyl)aniline,\n3-chloro-4-amino-N,N-bis-(2-hydroxyethyl)aniline,\n4-amino-N-ethyl-(piperidonoethyl)aniline,\n3-methyl-4-amino-N-ethyl-xcex2-(piperidonoethyl)aniline,\n4-amino-N-ethyl-N-(morpholinoethyl)aniline,\n4-amino-N-ethyl-N-(acetylaminoethyl)aniline,\n4-amino-N-(methoxyethyl)aniline,\n3-methyl-amino-N-ethyl-N-(2-acetylaminoethyl)aniline,\n4-amino-N-ethyl-N-(mesylaminoethyl)aniline,\n3-methyl-4-amino-N-ethyl-N-(xcex2-mesylaminoethyl)aniline,\n4-amino-N-ethyl-N-(xcex2-sulfoethyl)aniline,\n3-methyl-4-amino-N-ethyl-N-(xcex2-sulfoethyl)aniline,\nN-(4-aminophenyl)morpholine,\nN-(4-aminophenyl)piperidine,\n2,3-dimethyl-p-phenylenediamine,\n2-isopropyl-p-phenylenediamine,\nN,N-bis-(2-hydroxyethyl)-p-phenylenediamine sulfate,\no-aminophenol,\nm-aminophenol,\np-aminophenol,\n2-iodo-p-aminophenol,\n2-nitro-p-aminophenol,\n3,4-dihydroxyaniline,\n3,4-diaminophenol,\n2-hydroxy-4-aminobenzoic acid,\n2-hydroxy-4-aminobenzaldehyde,\n3-amino-4-hydroxybenzenesulfonic acid,\nN,N-diisopropyl-p-aminophenol,\nN-methyl-N-(1-propenyl)-aminophenol,\nN-phenyl-N-benzyl-p-aminophenol sulfate,\nN-methyl-N-(3-ethylphenyl)-p-aminophenol,\n2-nitro-5-ethyl-p-aminophenol,\n2-nitro-5-(2-bromoethyl)-p-aminophenol,\n(2-hydroxy-5-aminophenyl)acetaldehyde,\n2-methyl-p-aminophenol,\n(2-hydroxy-5-aminophenyl)acetic acid,\n3-(2-hydroxy-5-aminophenyl-1-propene,\n3-(2-hydroxy-5-aminophenyl)-2-chloro-1-propene,\n2-phenyl-p-aminophenol palmitate,\n2-(4-nitrophenyl)-p-aminophenol,\n2-benzyl-p-aminophenol,\n2-(4-chlorobenzyl-p-aminophenol perchlorate,\n2-(4-methylphenyl)-p-aminophenol,\n2-(2-amino-4-methylphenyl)-p-aminophenol,\np-methoxyaniline,\ndi-(2-aminoethyl-4-aminophenyl) ether,\ndi-(2-hydroxyethyl-4-aminophenyl) ether,\n(4-aminophenoxy)acetaldehyde,\n(4-aminophenoxy)acetic acid,\n(4-aminophenoxy)methanesulfonic acid,\n1-propenyl-4-aminophenyl ether isobutyrate,\ndi-(2-chloro-1-propenyl-4-aminophenyl) ether,\ndi-(2-nitro-1-propenyl-4-aminophenyl) ether,\ndi-(2-amino-propenyl-4-aminophenyl) ether,\ndi-(2-hydroxy-1-propenyl-4-aminophenyl) ether,\nN-methyl-p-aminophenol,\n3-methyl-4-aminophenol,\n2-chloro-4-aminophenol,\n3-chloro-4-aminophenol,\n2,6-dimethyl-4-aminophenol,\n3,5-dimethyl-4-aminophenol,\n2,3-dimethyl-4-aminophenol,\n2,5-dimethyl-4-aminophenol,\n2-hydroxymethyl-4-aminophenol,\n3-hydroxymethyl-4-aminophenol,\no-hydroxyphenol (catechol),\nm-hydroxyphenol (resorcinol),\np-hydroxyphenol (hydroquinone),\n4-methoxyphenol,\n2-methoxyphenol,\n4-(2-chloroethoxy)phenol,\n4-(2-propenoxy)phenol,\n4-(3-chloro-2-propenoxy)phenol,\n2-chloro-4-hydroxyphenol (2-chlorohydroquinone),\n2-nitro-hydroxyphenol (2-nitrohydroquinone),\n2-amino-4-hydroxyphenol,\n1,2,3-trihydroxybenzene (pyrogallol),\n2,4-dihydroxybenzaldehyde,\n3,4-dihydoxybenzoic acid,\n2,4-dihydroxybenzenesulfonic acid,\n3-ethyl-4-hydroxyphenol,\n3-(2-nitroethyl)-4-hydroxyphenol,\n3-(2-propenyl)-4-hydroxyphenol,\n3-(3-chloro-2-propenyl)-4-hydroxyphenol,\n2-phenyl-4-hydroxyphenol,\n2-(4-chlorophenyl)-4-hydroxyphenol,\n2-benzyl-4-hydroxyphenol,\n2-(2-nitrophenyl)-4-hydroxyphenol,\n2-(2-methylphenyl)-4-hydroxyphenol,\n2-(2-methyl-4-chlorophenyl)-4-hydroxyphenol,\n2-methoxy-4-(1-propenyl)phenol,\n4-hydroxy-3 methoxycinnamic acid,\n2,5-dimethoxyaniline,\n2-methylresorcinol,\naniline,\np-chloroaniline,\np-fluoroaniline,\np-nitroaniline,\np-aminobenzaldehyde,\np-aminobenzoic acid,\nsodium p-aminobenzoate,\nlithium p-aminobenzoate,\ncalcium di-(p-aminobenzoate),\nammonium p-aminobenzoate,\np-aminobenzenesulfonic acid,\npotassium p-aminobenzenesulfonate,\nN-methylaniline,\nN-propyl-N-phenylaniline,\nN-methyl-N-2-propenylaniline,\nN-benzylaniline,\nN-(2-ethylphenyl)aniline,\n4-methylaniline,\n4-(2-bromoethyl)aniline,\n2-(2-nitroethyl)aniline,\n4-aminophenylacetaldehyde,\n4-aminophenylacetic acid,\n4-(2-propenyl)aniline acetate,\n4-(3-bromo-2-propenyl)aniline,\n4-phenylaniline chloroacetate,\n4-(3-chlorophenyl)aniline,\n4-benzylaniline,\n4-(4-iodobenzyl)aniline,\n4-(3-ethylphenyl)aniline,\n4-(2-chloro-ethylphenyl)aniline,\nphenol,\np-chlorophenol,\np-nitrophenol,\np-hydroxybenzaldehyde,\np-hydroxybenzoic acid,\np-hydroxybenzenesulfonic acid,\nethylphenyl ether,\ndi-(2-chloroethylphenyl) ether,\ndi-(2-nitroethylphenyl) ether,\nphenoxyacetaldehyde,\nphenoxyacetic acid,\n3-phenoxy-1-propene,\n3-phenoxy-2-nitro-1-propene,\n3-phenoxy-2-bromo-1-propene,\n4-propyiphenol,\n4-(3-bromopropyl)phenol,\n2-(2-nitroethyl)phenol,\n4-hydroxyphenylacetaldehyde,\n4-hydroxyphenylacetic acid,\n4-(2-propenyl)phenol,\n4-phenyiphenol\n4-benzylphenol,\n4-(3-fluoro-2-propenyl)phenol,\n4-(4-chlorobenzyl)phenol,\n4-(3-ethylphenyl)phenol\n4-(2-chloro-3-ethylphenyl)phenol,\n2,5-xylenol,\n2,5-diaminopyridine,\n2-hydroxy-5-aminopyridine,\n2-amino-3-hydroxypyridine,\ntetraaminopyrimidine,\n1,2,4-trihydroxybenzene\n1,2,4-trihydroxy-5-(C1-C6-alkyl)benzene,\n1,2,3-trihydroxybenzene,\n4-aminoresorcinol,\n1,2-dihydroxybenzene,\n2-amino-1,4-dihydroxybenzene,\n2-amino-4-methoxyphenol,\n2,4-diaminophenol,\n3-methoxy-1,2-dihydroxybenzene,\n4,6-dimethoxy-3-amino-1-hydroxybenzene,\n2,6-dimethyl-4-(p-hydroxyphenyl)amino]phenol\nand salts thereof.\nAdditional oxidation dye couplers suitable for use herein include catechol species and in particular catechol xe2x80x9cdopaxe2x80x9d species which includes dopa itself as well as homologs, analogs and derivatives of DOPA. Other suitable dye precursors are dihydroxyindole (DHI), dihydroxyindolecarboxylic acid (DHICA) and derivatives thereof, indolines and derivatives thereof. Examples of suitable catechol species include cysteinyl dopa, alpha alkyl dopa having 1 to 4, preferably 1 to 2 carbon atoms in the alkyl group, epinephrine and dopa alkyl esters having 1 to 6, preferably 1 to 2 carbon atoms in the alkyl group.\nThe oxidation dye couplers (precursors) can be used herein alone or in combination with other oxidation dye couplers (precursors) mentioned above. The choice of a single dye coupler (precursor) will be determined by the color, shade and intensity of coloration which is desired. The following are preferred oxidation dye couplers (precursors) which can be used herein, singly or in combination, to provide oxidation hair dyes having a variety of shades ranging from ash blonde to black; these are: pyrogallol, resorcinol, p-toluenediamine, o-phenylenediamine, m-phenylenediamine, o-aminophenol, p-aminophenol, 4-amino-2-nitrophenol, nitro-p-phenylenediamine, N-phenyl-p-phenylenediamine, m-aminophenol, 2-amino-3-hydroxytoluene, N,N bis (2-hydroxyethyl)p-phenylenediamine, 4-amino-2-hydroxytoluene, 1,5-dihydroxynapthalene 2,4-diaminoanisole, hydroquinone, 4-amino-2-hydroxytoluene, 2-methyl resorcinol, 2-methyl-5-hydroxyaminophenol, 6-amino-3-hydroxy-toluene, 2,5-diaminotoluene, and 1-phenyl-3-methyl-pyazolone. These can be used in the molecular form or in the form of peroxide-compatible salts, as detailed above.\nApplicant has also found that compositions of the present invention comprising particular oxidizing agents with particular transition metal salts or catalysts with oxidative dyes are valuable for the delivery of good high intensity colors with reduced levels of dye and peroxide. In particular applicant has now found that good hair coloring results can be achieved using the oxidizing systems of the present invention and up to about 84% less dye versus conventional hair coloring compositions. As such the compositions according to the present invention are valuable for the delivery of good initial color and wash fastness over time in addition to having low hair damage attributes.\nSolvents\nWater is the preferred principal diluent for the compositions according to the present invention. As such, the compositions of present invention may also include one or more solvents as additional diluent materials. Generally, the solvent is selected to be miscible with water and innocuous to the skin. Solvents suitable for use herein include C1-C20mono- or polyhydric alcohols and their ethers, glycerine, with monohydric and dihydric alcohols and their ethers preferred. In these compounds, alcoholic residues containing 2 to 10 carbon atoms are preferred. Thus, a particularly preferred group includes ethanol, isopropanol, n-propanol, butanol, propylene glycol, ethylene glycol monoethyl ether, and mixtures thereof.\nThe compositions of the invention may also include the following materials.\nBuffering Agents\nThe coloring compositions of the present invention may have a pH in the range of from about 5 to about 11, more preferably from about 6 to about 8, and especially from about 6 to about 7.\nAs herein before described the preferred coloring compositions of the present invention may contain one or more buffering agents and/or hair swelling agents (HSAs) to adjust the pH to the desired level. Several different pH modifiers can be used to adjust the pH of the final composition or any constituent part thereof.\nThickeners\nThe composition containing one or more oxidative hair coloring agents of the present invention (coloring compositions) may additionally include a thickener at a level of from about 0.05% to about 20%, preferably from about 0.1% to about 10%, more preferably from about 0.5% to about 5% by weight. Thickening agents suitable for use in the compositions of the invention may be selected from the group consisting of oleic acid, cetyl alcohol, oleyl alcohol, sodium chloride, cetearyl alcohol, stearyl alcohol, and synthetic thickeners; and mixtures thereof.\nSurfactant Materials\nThe compositions of the present invention may additionally contain a surfactant system. Suitable surfactants for inclusion in the compositions of the invention generally have a lipophilic chain length of from about 8 to about 22 carbon atoms and can be selected from the group consisting of anionic, cationic, nonionic, amphoteric, zwitterionic surfactants and mixtures thereof.\nMaterials used in compositions of the invention may be employed in the following ranges:"}
-{"text": "The touch control technology used in touch panels can be classified into resistive type, capacitive type, surface acoustic wave type, and optics type, based on their sensing principles.\nFIGS. 1 and 2 are a block diagram and a circuit diagram of a control circuit of a conventional capacitive touch panel, respectively. The control circuit 10 is electrically connected to the capacitive touch panel 12, and comprises a pulse generator 20, a multiplexer 22, a multiplexer 24, an integrator 26, a sampling and holding circuit 28, and an analog to digital converter (hereinafter called the ADC) 32.\nThe pulse generated by the pulse generator 20 functions as a driving which is transmitted through the multiplexer 22 to the X-axial driving wires and the Y-axial driving wires of the capacitive touch panel 12, and the value of capacitive change sensed by the Y-axial and X-axial sensing wires of the capacitive touch panel 12 functions as a touch control sensing signal, which is transmitted through the multiplexer 24 to the integrator 26.\nIn FIG. 2, the touch control sensing signal of the capacitive touch panel 12 transmitted through the multiplexer 24 to the integrator 26 is a pulse signal, and the integrator 26 integrates several pulse signals into a ladder-shaped wave signal and transmits it to a sampling and holding circuit 28.\nThen, in FIG. 1, the sampling and holding circuit 28 samples the ladder-shaped wave signal. integrated by the integrator 26, and keeps the sampled signal at a direct current (DC) voltage level. The ADC 32 will convert the aforementioned sampled analog signal into a digital signal and transmit the digital signal to a microprocessor (not shown) to conduct a signal processing.\nThe aforementioned control circuit 10 used in the conventional capacitive touch panel 12, at least, has the following disadvantages:\nThere is noise signal generated in the circuit environment, and the noise signal will also be integrated by the integrator 26 as a portion of the signal, resulting in a low signal to noise ratio (SNR) of the whole system.\nIt takes a period of time for the integrator 26 to transfer and accumulate the charges on the capacitive touch panel 12 to a voltage that can be processed by a rear-end circuit. Such a long integration time will affect the frame rate of the capacitive touch panel 12. This problem may be compensated by way of a specific hardware (for example, the sampling and holding circuit 28); however, the cost will increase.\nThe use of the integrator 26, sampling and holding circuit 28, and digital to analog converter 32, will occupy a larger area on the chip."}
-{"text": "When a high-precision metal microstructure is mass-produced, a process of lithographie galvanoformung abformung (LIGA) is useful. The LIGA process, which uses synchrotron radiation (SR) ray among X-rays for its high directivity, is characterized in that it allows deep lithography, processing of a structure having a height of several hundreds of micrometers with a precision on the order of micrometers, easy manufacturing of a metal microstructure having high aspect ratio, and the like. Therefore, the LIGA process is expected to find application in various fields.\nThe LIGA process is a processing technique combining lithography, plating such as electroforming, and molding. According to the LIGA process, a resist film is formed on a conductive substrate, for example, and then is irradiated with SR ray through an absorbent mask having a pattern of predetermined shape to perform lithography. Such lithography allows a resist structure (resin mold) to be formed corresponding to the pattern shape of the absorbent mask. A vacant portion of the resist structure is then electroformed to obtain a metal microstructure. Furthermore, a high-precision metal microstructure obtained by further electroforming may be used as a mold to perform molding such as injection molding so as to obtain a micromolded article made of resin. For example, Japanese Patent Laying-Open No. 2002-217461 discloses a method of manufacturing a metal microstructure as a mold by the LIGA process to form a micromolded article.\nHowever, when lithography is used to form a metal microstructure, it can produce only one resist structure (resin mold) at a time. Therefore, lithography is not suitable for mass production. Furthermore, if a resin mold to be electroformed to form a metal microstructure is not fixed firmly with high precision onto a conductive substrate with an adhesive or the like, dimensional change or separation between the resin mold and the conductive substrate occurs in a subsequent step, disadvantageously resulting in a product of lower precision, a product not suitable for commercial use, and the like. Furthermore, in the case where a typical adhesive is applied to the conductive substrate to fix the resin mold thereon so as to solve the problem of a product of lower precision and a defective product, if the adhesive remains at a vacant portion of the layered structure having the resin mold, electromolding cannot be performed or insufficiently performed thereon, disadvantageously resulting in an uneven metal layer formed thereby. Furthermore, if the entire layered structure having the resin mold is etched to remove the adhesive remaining at a vacant portion thereof, the resin mold itself is damaged accordingly to a certain degree.\nAn object of the present invention is to provide a method of manufacturing a metal microstructure by using a resin mold, in which a mild condition can be set so as to cause less damage to a resin mold and a high-precision metal microstructure can be mass-produced by uniform electroforming."}
-{"text": "1. Field of the Invention\nThe present invention relates to a device and a method for transposing a matrix of video signals, which may be applied to a video decoder for decoding a matrix of compressed video bitstream such as signals of an MPEG(Moving Picture Expert Group) format.\n2. Discussion of the Related Art\nAs is well known, recently high definition television(HDTV) broadcasting systems have been developed, some of which are put into pilot operation in some of countries. In the pilot operation of the HDTV broadcasting systems, signals according to the HDTV broadcasting systems are transmitted, which are capable of being processed in an HDTV receiver (hereinafter called as HDTV). In the meantime, there are two classes of signals in the HDTV system of the ATSC (Advanced Television Systems Committee) of the U.S.A. One class of signal is an HD class signal for producing images of a high definition, and the other class of signal is an SD (Standard Definition) class signal(hereinafter called as SD signal) for producing images of a standard definition. The SD class signal has a data transmission rate lower than the HD class signal. The HD TV is adapted to process the two classes of video data, basically. Therefore, there can be four cases of signal transmission between broadcasting stations and HD TV\"\"s; the first case is displaying a HD signal as an HD signal, the second case is displaying an HD signal as an SD signal, the third case is displaying an SD signal as HD class signal, and the fourth is displaying an SD signal as an SD signal. Though the first, third and fourth cases have no particular problems in processing the signals with regard to the system performances, the second case does. That is, the second case has problem in that the images appear broken. Therefore, if an HD signal is to be displayed as an SD class signal, it is necessary to degrade performance of a received HD signal to some extent. For example, an HD signal of 20 MHz MPEG format should be altered to an SD class signal of 6 MHz MPEG format. This degradation of the HD signal performance implies a drop in the resolution. In general, the HD signal performance is degraded by removing a portion of frequency band over a certain frequency from the HD signal. An HD signal of an 8xc3x978 matrix having a portion of frequency band over a certain frequency removed therefrom to suit to the SD class signal is thereby altered into an 8xc3x974 HD signal. This 8xc3x974 HD signal can be displayed as an SD class signal, with a resolution degraded compared to the 8xc3x978 HD signal, but without the problem of broken images. In the meantime, an 8xc3x978 matrix SD signal can be displayed on an HDTV as it is. Because of the aforementioned reason, the HD TV should be adapted to process the 8xc3x978 HD signal as well as the 8xc3x974 HD signal. A general HD TV will be explained with reference to FIGS. 1xcx9c3.\nReferring to FIG. 1, upon receiving a compressed video bit stream from a broadcasting station, the general HD TV synchronizes a desired video bitstream in a tuner 10, of the video bitstream received through an antenna Ant, and demodulates in a demodulator 20. This video bitstream, being an MPEG signal, has a form of matrix. Then, the demodulated video bit stream is restored in a video decoder 30 and processed to a displayable condition in a VDP(Video Display Processor) 40. The signal from the VDP 40 is displayed through a displayer 50. In this instance, the VDP 40 processes the decoded signals appropriate to a performance of the displayer 50.\nFIG. 2 illustrates a detail block diagram of the video decoder 30 shown in FIG. 1. Upon reception of a matrix of the demodulated video bit stream from the demodulator 20 in FIG. 1, a VLD(Variable Length Decoder) 21 therein decodes the demodulated video bit stream to provide DCT(Discrete Cosine Transform) coefficients and motion vectors. The DCT coefficients are scanned by an inverse scanner 22 and inverse quantized by the inverse quantizer 23. Then, an IDCT(Inverse Discrete Cosine Transformer) 25 makes an inverse discrete cosine transformation of the inverse quantized DCT coefficients to provide spatial pixel values. Here, before the inverse discrete cosine transformation in the IDCT 25, the inverse quantized DCT coefficients are transposed in a transposer 24 for easy inverse discrete cosine transformation. On the other hand, a motion compensator 27 uses the motion vectors from the VLD 21 in compensating a reference video frame stored in the frame memory in advance, and an adder 28 adds a signal from the motion compensator 27 and a signal from the IDCT 25, to provide an added value to the VDP 40 in FIG. 1.\nIllustrated in FIG. 3 is an example of the transposer 24 for making alternative read/write of two memories 24a and 24b in transposing them.\nReferring to FIG. 3, the transposer 24 is provided with two SRAM(Static Random Access Memory) 24a and 24b each of a 16 bitsxc3x9764 words size and two multiplexed flipflops 24d and 24e. Upon reception of an enable signal in a read/write controller 24c, the read/write controller 24c provides a read/write controlling signal and a selecting signal to the first SRAM 24a and the second SRAM 24b. . The first SRAM 24a and the second SRAM 24b are operative in opposite manner in response to the read/write controlling signal. That is, a first data is written on the first SRAM 24a at the first time, the second SRAM 24b is left at a disabled state in response to the selecting signal. Then, while the first SRAM 24a is read, the second a SRAM 24b is written of a second data. Thus, the first SRAM 24a and the second SRAM 24b are read and written alternatively. The alternative first, second data from the two SRAMs 24a and 24b are transposed by the multiplexed flipflops 24d and 24e. Besides, there are a transposing method disclosed in a U.S Pat. No. 4,769,790 by time delaying and a U.S. Pat. No. 5,418,487 using dual port memories. These two Patents have a disadvantages either in a size of hardware or in a complicated hardware. And, as has been explained, the background art transposer 24 shown in FIG. 2 transposes an mxc3x97m matrix signal(for example, m=8) using two memories of mxc3x97m words capacity each. Because of the use of two memories, such a background art transposer has disadvantages in that the inverse discrete cosine transformer and the video decoder 30 shown in FIG. 2 have increased areas and comparatively low operation speed of 100 M sample/sec.\nAccordingly, the present invention is directed to device and method for transposing a matrix of video signals and a television receiver employing the same that substantially obviates one or more of the problems due to limitations and disadvantages of the related art.\nAn object of the present invention is to provide a method and a device for transposing a matrix of video signals which can reduce a chip area of the inverse discrete cosine transformer and, further, an area of the video decoder.\nAnother object of the present invention is to provide a method and a device for transposing a matrix of video signals which is operable at fast speed.\nOther object of the present invention is to provide a video decoder and a television receiver employing the same method and device.\nAdditional features and advantages of the invention will be set forth in the description which follows, and in part will be apparent from the description, or may be learned by practice of the invention. The objectives and other advantages of the invention will be realized and attained by the structure particularly pointed out in the written description and claims hereof as well as the appended drawings.\nTo achieve these and other advantages and in accordance with the purpose of the present invention, as embodied and broadly described, there is provided a memory part for writing a matrix of video signals thereon, a write control circuit shift and writes rows of the video signals on the memory part by any one unit either of a row unit or a column unit. The video signals written on the memory are read from the memory by one unit different from the unit in the writing either of the row unit and the column unit by a read control circuit for transposing. Rows of video signals received at the next time are shifted, and written on portions of the memory part emptied in sequence due to the shift in the reading by one unit opposite to the unit in the reading either of the row unit and the column unit in sequence. The memory part may includes either one memory on which the matrix of video signals may be written or two memory of less capacity, resulting to a reduction of an occupied area of the memory part. Because rows of video signals received at the next time is shifted, and written on portions of the memory part emptied in sequence in the reading, a speed of the transposing becomes faster.\nIt is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory and are intended to provide further explanation of the invention as claimed."}
-{"text": "Security systems have traditionally been utilized for detecting, removing, etc. unwanted data (e.g. malware, spyware, unsolicited data, etc.). Oftentimes, it is advantageous for separate security systems to communicate, such as for sharing information regarding known, detected, etc. unwanted data. Unfortunately, traditional security systems have exhibited various limitations with respect to communication capabilities thereof.\nJust by way of example, traditional security systems are generally periodically updated and upgraded, and the frequency of such updates and upgrades is increasing such that new types of systems are frequently becoming available. While new security systems may sometimes be capable of being programmed to communicate with old security systems via protocols existing in the old security systems, requiring such programming is restrictive since the protocols utilized must be defined at the time when the old security system was developed. In addition, when the old security systems need a way to communicate with newer security systems, communication therebetween is impossible since the old security systems are not already equipped with the communication protocols used by the newer security systems.\nA similar problem exists for two security systems developed independently (e.g. as a result of different developers or, simply, due to a lack of internal communication between such developers). Both security systems may be released, but to pass messages from one another both products may be required to be re-released, which may be very costly. There is thus a need for addressing these and/or other issues associated with the prior art."}
-{"text": "For various apparel articles such as jackets, baseball caps, and other products, as well as accessory articles such as backpacks and cinch bags, it may be desirable to have adornments, such as a patch, attached thereto. In certain applications use of a removable adornment (e.g. a removable patch) is advantageous as the indicia on the apparel or accessory can be changed by the user by simply replacing the adornment.\nA removable adornment can be attached to apparel and other products and accessories (articles) via numerous fastening techniques. For example, a patch can be removably attached to an article of manufacture using Velcro\u00ae (a hook-and-loop fastener or attachment system), magnets, adhesives, snaps, buttons, and pins. However, such attachment techniques often result in the patch protruding away from the surface of the apparel or product, which can look low-cost, unprofessional, and/or permit the patch to be unintentionally or accidently removed. In addition, the entire surface of the patch including its borders may not be flush with the surface of the apparel or accessory thereby causing wrinkles or other unevenness of the patch and its borders, again resulting in a low-cost, unsophisticated look.\nIn certain designs, the patch can have a border such as a stitched border around the patch. The stitched border can be a merrowed border or a satin stitched border. Although such borders on a patch can improve the look of the patch on the apparel or accessory, similar issues can be present for removable patches with borders as with removable patches without borders (e.g., patches having hot cut borders).\nThus, there is a need to improve the removable attachment of adornments to apparel and other products and accessories, which attachment can include patch receptacles for removable patches, such that imperfections of the adornment cannot be readily seen and/or the resulting combination of the adornment and the article can present the look of a bordered adornment on the apparel or accessory, despite the adornment being borderless."}
-{"text": "This invention relates to a single frequency multi-transmitter telemetry system in which the transmitted data is retained when only one transmitter is transmitting and is discarded when more than one transmitter is transmitting.\nResearch programs in general aviation require the instrumentation of aircraft for test flight purposes. Space, weight, and power consumption are critical factors to be considered in the design of the system and instrumentation of the aircraft. Installation rapidness is also a design consideration in order to make the system cost-effective; primarily because of manpower and aircraft leasing costs.\nIt is the primary object of this invention to provide a telemetry system that will improve the test flight data availability, accelerate the installation process, and reduce data cost.\nAnother object of this invention is to provide a simple single frequency telemetry system that will deliver a substantial amount of data at low cost.\nOther objects and advantages of this invention will be apparent hereinafter in the specification and drawings."}
-{"text": "1. Field of the Invention\nThe present invention generally relates to a printed circuit board. More particularly, the present invention relates to a printed circuit board integrated with a two-axis fluxgate sensor and a method for manufacturing the same.\n2. Description of the Related Art\nThe existence of magnetic energy has been proven through various physical phenomena. For example, a fluxgate sensor enables a human to indirectly perceive magnetic energy, since magnetic energy is not perceivable by human sense organs, such as eyes or ears. Regarding a fluxgate sensor, a magnetic sensor employing a soft magnetic coil has been conventionally used. A magnetic sensor is made by winding a coil around a relatively large bar-shaped core or annular-shaped core, which is formed of a soft magnetic ribbon. In addition, an electronic circuit is employed to obtain a magnetic field in proportion to the measured magnetic field.\nThe conventional fluxgate sensor, however, has the following problems, some of which are due to the structure of the conventional fluxgate sensor. In the structure of a conventional fluxgate sensor the coil being wound around the large bar-shaped core or annular-shaped core, which is made of the soft magnetic ribbon, results in high production costs and a large volume of the overall system.\nIn addition, flux leakage is inevitable in the flux change due to the excitation coil and the detected magnetic field. Accordingly, high sensitivity is not readily achieved.\nFurther, a conventional fluxgate sensor for detecting magnetic energy in two axes includes two separate one-axis fluxgate sensors arranged in perpendicular relation with each other, which results in yet another problem of the device size increasing."}
-{"text": "1. Field\nExemplary embodiments of the present invention relate to a semiconductor device fabrication technology, and more particularly, to an image sensor, and a method for fabricating the image sensor.\n2. Description of the Related Art\nAn image sensor is a device that transforms optical images of one or more dimensions on a photographic subject into electric signals. An image sensor includes a pixel array region for receiving incident light and generating photocharges and a logic region for generating electric signals by processing the generated photocharges. The logic region may also be called a peripheral circuit region. Incident light radiates not only upon the pixel array region but also upon the logic region as well. The photocharges generated in the logic region by the incident light radiating thereon may act as noise on the electric signals generated in the logic region, that is, electric signals generated by processing the photocharges generated in the pixel array region. Therefore, the logic region may use a light shielding pattern for blocking out the incident light.\nGenerally, a light shielding pattern is formed over a substrate of the logic region. Therefore, there is a step height difference between the pixel array region and the logic region. The step height difference between the two regions that is caused by the light shielding pattern deteriorates the characteristics and integration of the image sensor."}
-{"text": "As a method for manufacturing a semiconductor device, a method has been conventionally known of sealing one or more semiconductor chips fixed to a substrate, etc. with a sealing resin and dicing the sealed body to form a packaged semiconductor device unit. For example, a sealing sheet constituted with a thermosetting resin has been known as the sealing resin (for example, refer to Patent Document 1)."}
-{"text": "This invention relates to certain herbicidal ethers, agriculturally suitable compositions thereof, and a method for their use as broad spectrum preemergent or postemergent herbicides.\nNew compounds effective for controlling the growth of undesired vegetation are in constant demand. In the most common situation, such compounds are sought to selectively control the growth of weeds in useful crops such as cotton, rice, corn, wheat and soybeans, to name a few. Unchecked weed growth in such crops can cause significant losses, reducing profit to the farmer and increasing costs to the consumer. In other situations, herbicides are desired which will control all plant growth. Examples of areas in which complete control of all vegetation is desired are areas around railroad tracks, storage tanks and industrial storage areas. There are many products commercially available for these purposes, but the search continues for products which are more effective, less costly and environmentally safe.\nPayne et al., U.S. Pat. No. 4,567,283 and Payne et al., U.S. Pat. No. 4,670,041 disclose a variety of herbicidal bicyclic ethers of the Formula ##STR1##\nU.S. Pat. No. 4,798,621 and WO 8,902,219 both disclose bicyclic ethers and their method-of-use in rice.\nU.S. Pat. No. 4,486,219 discloses bicyclic ethers of the formula: ##STR2##\nWO91/03464 discloses bicyclic ethers such as: ##STR3## wherein inter alia X and Y are independently C(R.sup.3) (R.sup.4)OR.sup.5 ;\nR.sup.5 is C.sub.1 -C.sub.3 alkyl, C.sub.2 -C.sub.4 alkenyl, C3-C.sub.4 alkynyl, CH.sub.2 -phenyl or C.sub.1 -C.sub.4 alkyl substituted with OR.sup.8 or OCF3; PA1 R.sup.2 is H, C.sub.1 -C.sub.4 alkyl, C.sub.2 -C.sub.4 alkenyl or C.sub.2 -C.sub.4 alkynyl; PA1 R.sup.3 and R.sup.4 are independently H, C.sub.1 -C.sub.4 alkyl, C.sub.2 -C.sub.4 alkenyl, C.sub.2 -C.sub.4 alkynyl or C.sub.1 -C.sub.3 alkyl substituted with OCH.sub.3 or OCH.sub.2 CH.sub.3 ; ##STR5## X is CR.sup.5 R.sup.6 Y.sup.1, CR.sup.5 .dbd.CR.sup.6 Y.sup.3 or R.sup.5 R.sup.6 PA1 R.sup.5 and R.sup.6 are independently H or C.sub.1 -C.sub.3 alkyl PA1 Y.sup.1 is H; C.sub.1 -C.sub.4 alkyl optionally substituted with Y.sup.2 ; C.sub.3 -C.sub.6 cycloalkyl; CN; C(O)NR.sup.7 R.sup.8 ; C(O) NHOR.sup.9 ; CO.sub.2 R.sup.10 ; C(O)R.sup.11 C(OR.sup.22) (OR.sup.23)R.sup.11 ; CHO; CH(OR.sup.22) (OR.sup.23); CH-NOR.sup.16 CR.sup.11 .dbd.NOR.sup.17 ; SR.sup.15 ; halogen; C.sub.2 -C.sub.4 alkenyl optionally substituted with 1-3 halogens, CO.sub.2 (C.sub.1 -C.sub.3 alkyl), CN or CHO; C.sub.2 -C.sub.4 alkynyl optionally substituted with halogen; S(O).sub.n R.sup.12 ; SO.sub.2 NR.sup.13 R.sup.14 ; NO.sub.2 N.sub.3 SCN; phenyl optionally substituted with 1-3 substituents selected from CH.sub.3, OCH.sub.3, SCH.sub.3, CF.sub.3 and halogen; imidazole; triazole; tetrazole; or ##STR6## Y.sup.2 is CN; C(O)NR.sup.7 R.sup.8 ; C(O)NHOR.sup.9 ; CO.sub.2 R.sup.10 ; C(O)R.sup.11 ; CHO; CH.dbd.NOR.sup.16 ; CR.sup.11 .dbd.NOR.sup.17 ; SR.sup.15 ; 1-3 halogens C.sub.2 -C.sub.4 alkenyl optionally substituted with 1-3 halogens; C.sub.2 -C.sub.4 alkynyl optionally substituted with halogen; S(O).sub.n R.sup.12 ; SO.sub.2 NR.sup.13 R.sup.14 ; NO.sub.2 ; N.sub.3 ; SCN; phenyl optionally substituted with 1-3 substituents selected from CH.sub.3, OCH.sub.3, SCH.sub.3, CF.sub.3 and halogen; imidazole triazole tetrazole or ##STR7## Y.sup.3 is H, C.sub.1 -C.sub.3 alkyl, CN, CO.sub.2 R.sup.10, C(O)R.sup.11, CHO or halogen; PA1 R.sup.7 R.sup.8 R.sup.13 and R.sup.14 are independently H or C.sub.1 -C.sub.3 alkyl; R.sup.7 and R.sup.8 may be taken together to form a 5- or 6-membered ring and R.sup.13 and R.sup.14 may be taken together to form a 5- or 6-membered ring; PA1 R.sup.9, R.sup.12, R.sup.16 and R.sup.17 are independently C.sub.1 -C.sub.3 alkyl; PA1 R.sup.10 is H; C.sub.1 -C.sub.6 alkyl optionally substituted with 1-3 halogens, OR.sup.18, CN or phenyl optionally substituted with 1-3 substituents selected from halogen, CH.sub.3 and OCH.sub.3 ; C.sub.3 -C4alkenyl; or C.sub.3 -C.sub.4 alkynyl; PA1 R.sup.11 is C.sub.1 -C.sub.6 alkyl optionally substituted with 1-3 halogens, OR.sup.18, CN or phenyl optionally substituted with 1-3 substituents selected from halogen, CH.sub.3 and OCH.sub.3 ; C.sub.2 -C.sub.4 alkenyl; C.sub.2 -C.sub.4 alkynyl; or phenyl optionally substituted with 1-3 substituents selected from halogen, CH.sub.3 and OCH.sub.3 ; PA1 R.sup.15 is H, C.sub.1 -C.sub.3 alkyl or benzyl; PA1 n is 1 or 2; PA1 Q is CH.sub.2 W or ##STR8## q and r are independently 0-2; R.sup.18 is H or C.sub.1 -C.sub.3 alkyl; PA1 R.sup.19 is H, halogen, C.sub.1 -C.sub.3 alkyl, OR.sup.20, SR.sup.20 or CN; PA1 R.sup.20 is C.sub.1 -C.sub.3 alkyl or C.sub.1 -C.sub.3 haloalkyl; PA1 Z is CH.sub.2, NR.sup.21, O, S or may be CH and taken to form a double bond with an adjacent carbon; PA1 R.sup.21 is H or C.sub.1 -C.sub.3 alkyl; PA1 R.sup.22 and R.sup.23 are independently C.sub.1 -C.sub.3 alkyl or may be taken together to from a 5- or 6-membered ring; PA1 R.sup.24, R.sup.25 and R.sup.26 are independently H or C.sub.1 -C.sub.2 alkyl; PA1 W is phenyl optionally substituted with 1-3 substituents selected from halogen, C.sub.1 -C.sub.3 alkyl, C.sub.1 -C.sub.3 alkoxy, OH, CN, C.sub.1 -C.sub.3 haloalkyl, C.sub.1 -C.sub.3 haloalkoxy, C.sub.1 -C.sub.3 alkylthio, C.sub.2 -C.sub.4 alkenyl and C.sub.2 -C.sub.4 alkynyl or W is a 5-, 6- or 7-membered heterocyclic ring containing 1 to 3 heteroatoms selected from the group 0-2 nitrogens, 0-2 oxygens and 0-2 sulfurs, each ring optionally substituted with 1-2 substituents selected from halogen, CH.sub.3 and OCH.sub.3 ; PA1 1) in compounds of Formulas I and II, Y.sup.1 is other than H, C.sub.1 -C.sub.4 alkyl or halogen, and Y.sup.2 is other than halogen(s); PA1 2) in compounds of Formula IV, when Q is other than CH.sub.2 W, then Y.sup.1 is other than H or C.sub.1 -C.sub.4 alkyl; PA1 3) the sum of q and r is 0-2; and PA1 4) if the sum of q and r is O then Z is CH.sub.2. PA1 1. Compounds of Formulae I, II, III or IV wherein: PA1 3. Compounds of Preferred 2 wherein: PA1 4. Compounds of Preferred 3 wherein: PA1 6. Compound of Preferred 5 wherein: PA1 7. Compounds of Preferred 6 wherein the compound is of Formula I. PA1 8. Compounds of Preferred 6 wherein the compound is of Formula II. PA1 9. Compounds of Preferred 6 wherein the compound is of Formula III. PA1 10. Compounds of Preferred 6 wherein the compound is of Formula IV. PA1 (2-endo, 4-endo )-(.+-.)-4-[(2,6-difluorophenyl)-methoxy]-5-methyl -6-oxa-bicyclo[3.2.1 ]octane-2-acetonitrile; (2-endo, 4-endo)-(.+-.)-5-ethyl-4-[(2-fluorophenyl) -methoxy ]-6-oxabicyclo[3.2.1]octane-2-acetonitrile PA1 (2-endo, 4-endo )-(.+-.)-5-ethyl-4-(phenylmethoxy)-6-oxabicyclo [3.2.1 ]octane-2-acetonitrile (2-endo, 4-endo)-(.+-.)-4-[(2,6-difluorophenyl)-methoxy]-5-ethyl-6-oxabicyclo [3.2.1 ]octane-2-acetonitrile PA1 (2-endo, 4-endo)-(.+-.)-5-ethyl-4-[(2-fluorophenyl)-methoxy]-5-methyl-6-oxabicyclo[ 3.2.1]octane-2 -acetonitrile; PA1 (2-endo, 4-endo)-(.+-.)-4-[(2-fluorophenyl)methoxy]-5-methyl-6-oxabicyclo[3.2.1]oct ane-2-acetonitrile; PA1 (2-endo, 4-endo )-(.+-.)-2-(bromomethyl)-5-methyl-4-(phenylmethoxy)-6-oxabicyclo[3.2.1]oct ane; PA1 (2-endo, 4-endo)-(.+-.)-5-ethyl-4-[(2-fluorophenyl) -methoxy]-2-(iodomethyl)-6-oxabicyclo[3.2.1]-octane; PA1 (2-endo, 4-endo)-(.+-.)-4-[(2-chlorophenyl)methoxy]-5-methyl-6-oxabicyclo[3.2.1]oct ane-2-acetonitrile; PA1 (2-endo, 4-endo)-(.+-.)-2-(iodomethyl)-5-methyl-4-(phenylmethoxy)-6-oxabicyclo[3.2. 1]octane; and PA1 (2-endo, 4-endo)-(.+-.)-5-ethyl-2-[(methylthio)methyl]-4-(phenylmethoxy)-6-oxabicyc lo[3.2.1]octane.\nIt is an object of the present invention to provide compounds and compositions which exhibit herbicidal activity on a variety of weed species. It is a further object of the present invention to provide compounds and compositions that are herbicidally safe to rice, cereals and broadleaf crops. It is a feature of the present invention to furnish novel oxabicyclo ethers that exhibit useful herbicidal activity. These and other objects, features and advantages will become apparent with respect to the following description of the invention."}
-{"text": "A sugarcane harvesting machine of this type is known from U.S. Pat. No. 3,830,046. An inclined conveyor is situated at the end of the product flow through this machine, and conveys sugarcane sections, formed by comminuting the cut sugarcane stalks in the machine, into a collection car. The inclined conveyor has two continuous traction mechanisms, designed as chains, which extend in parallel to one another and are connected by steel carriers. The chains run over upper drive wheels, and lower deflection wheels which are carried along and situated in the feed area of the sugarcane sections. The inclined conveyor is designed as a scraper conveyor; i.e., it has a robust sheet steel floor which forms a running surface over the entire length of the inclined conveyor and scrapes the carriers, thus upwardly conveying the sugarcane sections resting on the sheet steel floor.\nA first and a second extractor which are used to remove waste and the like from the sugarcane sections are situated in the sugarcane harvesting machine. The extractors operate via a rotor which generates an air flow which carries along the waste and conveys it to the outside. A first extractor is situated in a transition zone of the sugarcane sections on the inclined conveyor, while the second extractor is situated at the end of the inclined conveyor and removes waste from the sugarcane sections when they drop down at the upper end of the inclined conveyor and fall into the collection car.\nSuch inclined conveyors described above with reference to the prior art have many drawbacks. The steel chains wear out very quickly due to adhering sand and dust. In addition, they fall out of the division after even a few hours of operation due to the acting tension, causing the teeth to come out of the drive wheels. For this reason, the conveyor requires continuous monitoring and maintenance, and in spite of this has a short service life. A high level of wear and high energy consumption occur during operation of the inclined conveyor, since the sugarcane sections are scraped during conveyance over the steel floor of the conveyor. This is also associated with a high noise level due to the metal-to-metal contact. Furthermore, the scraper floor of the conveyor, which is made of solid sheet steel, is very heavy. A further drawback is that, due to the scraper floor which acts as a barrier, the inclined conveyor, although it is perforated, conveys almost all of the impurities along with it, resulting in only a very low screening and cleaning effect and thus subjecting the second extractor to heavy load."}
-{"text": "Capacitance material level probes of the type to which the present invention is directed typically comprise a central rod surrounded by a first layer of insulation, with the rod tip being exposed to act as a measurement antenna when mounted within a storage vessel. A tubular guard shield surrounds the first insulation layer and is itself surrounded by a second layer of insulation. A portion of the tubular guard shield adjacent to the rod tip is exposed and is adapted to be electrically energized separately from the rod so as to direct probe energy outwardly into the tank rather than through any material which may be coated onto the probe body. The second insulation layer is clamped within a nipple which is adapted to be threaded into a gland on a material vessel wall with the probe rod and guard shield projecting into the tank as described.\nMeasurement probes of the described character heretofore proposed in the art typically include the insulation layers and the guard shield as separate tubular elements telescopically received over the central rod. To obtain desired sealing within a pressurized vessel, it has often been necessary to crimp or squeeze the tip of the tubular guard shield around the first insulation layer to close off a leakage path therebetween. The separate tubular insulation elements and assembly thereof to the probe rod are expensive, and it is now proposed to perform the assembly operation and provide the necessary insulation in a single step by injection molding insulation material surrounding and between fixtured guard shield and probe elements. However, considerable shrinkage is encountered with thermosetting plastics suitable for this application, and crimping of the guard shield tip, which is undesirable in any event because it requires an extra assembly operation, is not entirely effective.\nAn object of the present invention, therefore is to provide an improved capacitance-type probe of the character described which is adequately and securely sealed at the tip of the guard shield.\nAnother object of the invention is to provide an economical method of manufacturing such a probe."}
-{"text": "A current method for the large scale collection and concentration of solar energy is the use of heliostat mirrors to reflect the sun's rays to a central receiver. By utilizing multiple heliostat mirrors, each one reflecting to a common point, concentration of solar energy is achieved.\nIn known systems, heliostat mirrors may be in a fixed position surrounding the tower. The mirror surfaces are typically controlled in two degrees of motion to position the surface of the mirror with respect to the tower. Each heliostat mirror has a control system which tracks the motion of the sun with respect to the centrally located receiver. The mirror is on a fixed base and the reflective surface of the mirror is continuously moved to maintain the solar reflection from the surface of the mirror onto the receiver. The purpose of positioning the heliostat mirrors being to reflect and direct the sun's rays to a designated central collection point, known as a central target receiver or a power tower. In order to accomplish this, the heliostat mirror requires a surface area of reflective mirror, two axes of motion, a servo motor for each axis of motion and a control system for positional calculation and motion control of the two axes.\nIn some installations, the array of heliostats encompasses a full 360 degree area around a receiver. This is a very expensive installation due to the large number of heliostats required in order to populate the array and surround the receiver. Also, because the position of the sun is changing throughout the day, at any point in time more than half of the array has a very low performance due to their fixed placement within the array.\nOther installations have attempted to address the issue of heliostat cost by reducing the number of mirrors. A reduced number of mirrors are controlled as a gang in that the heliostat array is a single movable unit that is positioned on a track. The array then rotates concentrically about the receiver on a horizontal plane of motion. The mirrors in the array are not only controlled as a single unit about the receiver, but are also movable as a single unit about a vertical axis to rotate their surface position with respect to the receiver. These installations are not practical for an installation having a large radius from the central tower. Also, large arc segments of mirrors will not target the receiver with adequate accuracy.\nFurthermore, known concentric designs do not always take into account the requirement of a full 360\u00b0 for rotation about the central receiver. In global implementations, particularly for solar tracking in latitudes between the Tropic of Cancer and the Tropic of Capricorn which comprises a large global market segment for solar thermal applications, known concentric designs may not capable of tracking the sun. For this special case, the solar path is not always southerly or northerly, but both depending on the day of the year. In fact, the solar path may be north or south of the central receiver and at one point, will be directly overhead. Without the capability of a full circle of motion for the array, known concentric designs will not operate within this equatorial band.\nAnother major factor in a solar thermal concentration system is the environmental impact of the installation. In known systems, a large area of land is necessary and that land must be dedicated to the sole purpose of solar thermal power collection. Additionally, the typical ground preparation for installation involves invasive grading and leveling practice which completely destroys the indigenous habitats and ecology of the terrain. Furthermore, the dedicated area is permanently covered by the mirror array, limiting ground exposure to natural heal and light provided by the sunlight, which negatively impacts the environment.\nThere is a need for a solar thermal concentration system and method that is cost effective to install and maintain, fully utilizes the collection possibilities as the earth orbits the sun, and has minimal adverse impact on the environment."}
-{"text": "1. Field of the Invention\nAn embodiment relates to a light emitting apparatus and a light unit.\n2. Description of the Related Art\nIn recent years, a light emitting apparatus and a lighting system using a light emitting device as a light source have been widely used. The light emitting diode (LED) may configure the light source that generate light by using compound semiconductor materials such as GaAs-series, AlGaAs-series, GaN-series, InGaN-series, and InGaAlP-series.\nMeanwhile, since the emission characteristics of the LED are deteriorated due to heat, the light emitting apparatus and the lighting system that have the LED need to be designed to rapidly emit heat emitted from the LED to the outside."}
-{"text": "Concerns about the environmental impact of fluorocarbons and hydrofluorocarbons on the environment call for expanded uses of environmentally sound refrigerants, such as water or ammonia. For commercial chiller applications in the capacity range of 15-10,000 refrigeration tons with little heating requirements, water is usually the preferred refrigerant due to its non-flammable and benign nature.\nAqueous absorption fluid cycles taking advantage of such refrigerants have been known and used for many decades. Single effect and various two-stage designs are commercially used in many countries. However, the increasing concern about the overall CO.sub.2 generation in the process of converting fossil fuels into energy used for air conditioning or refrigeration calls for higher energy conversion efficiencies than currently obtained with single stage (COP=0.6 to 0.8) or two-stage (COP=0.9 to 1.25) absorption equipment.\nU.S. Pat. No. 4,732,008 teaches the use of two single stage cycles coupled to achieve three refrigeration effects. The independent loops carry different absorption fluids with the lower stage using fluids such as aqueous LiBr solutions. However, the needed upper stage fluid crystallization and vapor pressure suppression properties require the use of a different fluid. Performance estimates lead to COPs of 1.5 to 1.7 if suitable upper stage fluids are used."}
-{"text": "Memory systems using \"hard disks\" to store information in annular tracks concentrically arranged in magnetic material formed on a disk surface are well known in the prior art. The information (often called \"data\") is stored in magnetic domains and each domain typically stores an encoded data bit. A major advantage of this type of memory system is that information can be easily stored on a selected track on a disk surface as the disk rotates at high speed and can be randomly accessed at high speed to store information or to read out information stored on a given track. A read-write head or transducer having a small coil and a magnetic core is moved across the disk surface to the track on which the desired information is to be stored, or from which the desired information is to be retrieved while the disk is spinning. The read-write head stores or detects data in a selected data track on the disk surface. Ideally, the head rides on a thin layer of air generated by the spinning disk. It is then said to fly over the disk surface. However, due to imperfections in the disk surface or to other operating factors such as system vibrations, the head may come into undesired contact with the magnetic film during flight while the disk is spinning, causing instant friction between the head and the magnetic film which may ultimately lead to the destruction of the film and to a head crash. Also, during start/stop there is intimate contact between the head and the disk surface with associated friction which may lead to excessive wear, head crash, or stiction. Stiction may be defined as head/disk adhesion strong enough to require manual removal of the head, resulting in severe marks on the disk surface. Stiction is often encountered when using relatively thick liquid lubricant or when the memory system is operating under conditions of high relative humidity.\nIn the prior art various techniques have been used to reduce the friction between the head and the magnetic film and the associated damage caused by it.\nIn memory systems which utilize \"flexible\" media such as magnetic tape or floppy disks, the space between the moving medium and the read/write head is much smaller than in hard disk systems, and there is continuous intimate contact between the medium and the read/write head. In both hard and flexible media, a magnetic film is formed on a substrate. For a hard disk, the substrate is typically an aluminum alloy with a relatively thick and hard undercoat formed thereon. The magnetic film is formed on the undercoat and is typically a cobalt phosphorus alloy or particulate oxide. The technique used to reduce friction or contact damage depends on the composition and quality of the magnetic film.\nOne method of forming a magnetic film on a disk substrate involves suspending iron-oxide power in a liquid media containing resins and coating the disk surface with the liquid. The disk is then cured in a strong magnetic field to provide the desired anisotrophy of magnetization. Such particulate magnetic films are typically formed to a thickness of between 0.7 and 5.0 microns, although recently particulate films having a thickness approaching 0.5 micron have been employed. These particulate iron-oxide magnetic films have a porous surface, and for these surfaces, one antifriction coating technique consists in applying a thin film of an organic lubricant, i.e., an organic liquid or wax such as Fomblin or Bareco. These organic lubricants are not very satisfactory because they tend to spin off the disk when the disk is rotated at operating speeds, e.g., 3600 rpm.\nThe trend in magnetic films has been to thinner films, because generally speaking, the thinner the magnetic film, the greater the information storage density, and to metal alloy film, as opposed to particulate oxide films. Particulate oxide films are porous and have low remanance due to their relatively small volume fraction of magnetic material in the resin. Since the magnitude of the read back signal is directly related to the remanance of the film, particlate oxide films are necessarily thick relative to alloy films of the same remanence.\nThe quest for thin magnetic films has led to several deposition techniques. Thin films typically comprise magnetic alloys of cobalt with such elements as phosphorus, nickel, chromium, or tungsten. Methods of deposition of thin magnetic alloy films include electroless, electroplating, sputtering, evaporation and other vapor deposition techniques. Thin magnetic films may also be formed by vapor deposition or sputtering or a magnetic metal oxide such as iron oxide in the form of Fe.sub.3 O.sub.4 or Fe.sub.2 O.sub.3 or by oxidation of metallic iron or by thermal decomposition of iron salts. Thin magnetic films typically have a thickness of 600-2000 .ANG., although even thinner films are extant, some having a thickness in the range of 300 .ANG.. These thin magnetic films generally have a smooth, nonporous surface, and the prior art liquid or waxy lubricants that were marginally satisfactory for particulate films, are less satisfactory for smooth, thin films since a liquid or waxy lubricant spins off the disk after a short period of operation.\nPowdered non-metallic solid lubricants, such as graphite and molybdenum disulfide have been applied to plated thin magnetic films by rubbing under pressure, but have proven to be unsatisfactory. (See Yanagisawa, U.S. Pat. No. 4,398,562, issued June 28, 1983). In general, these lubricants are made of loose particles which often detach and break, leading to a head crash. The particles are too large, e.g. 3 microns, average diameter, for an adequate thin layer of lubricant, and they present poor adhesion to metallic magnetic films.\nSince no satisfactory antifriction lubricants have been devised for thin magnetic films, the trend has been to use an extremely hard protective coating, such as plated rhodium and nickel-tin, which have a hardness in the range of 500-1000 Kg/mm.sup.2, and sputtered carbon, which has a hardness in the range of 3,000-5,000 Kg/mm.sup.2, as shown by R. J. Gambino and J. A. Thompson. Solid State Communications, Vol. 34 pp. 15-18 (1980). While these hard coatings tend to protect the magnetic medium from the head riding above the disk they have several drawbacks: one, they are generally expensive to implement; two, since they are so extremely hard, they tend to cause wear of the head riding on their surface. This tends merely to shift the wear problem from the disk to the head (and increasing the hardness of the head tends to shift the problem back to the disk).\nIn the case of sputtered carbon, the sputter deposition equipment is every expensive, the throughput low, and it is difficult to control its adhesion to the magnetic film and its crystalline character.\nRhodium, which is often taught, is still unsatisfactory in its performance. To assure effective protection, a thick layer of rhodium, typically on the order of about ten micro inches, has been employed. See U.S. Pat. No. 3,767,369, \"Duplex Metallic Overcoating\", issued Oct. 23, 1973, which is incorporated herein by reference. The U.S. Pat. No. 3,767,369 proposes an alternative to rhodium comprising a duplex of an at least 2 micro inches thick layer of an alloy of nickel and tin overlaid by a layer of rhodium at least 2 micro inches thick. The nickel-tin layer contains a proportion of nickel of about 35% by weight. Again, this nickel-tin alloy undercoat has a Vickers hardness of 500 to 900. More recently, a layer of rhodium 750 .ANG. thick has been suggested as a protective layer in combination with a thin layer of liquid lubricant, Dupont 804. See Rossi, et al, \"Vacuum Deposited Thin Metal Film Disk\", J. Appl. Phys., 55 (No. 6) 2254 (1984), which is incorporated herein by reference. Such a combination ignores the long-term reliability problem due to the spin-off of the lubricant. Also, the total thickness of the overcoat, including chromium, rhodium, and liquid lubricant is about 1,000 .ANG.."}
-{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display (LCD) device, and more particularly, to an LCD device that can reduce a fabrication cost and can prevent a decrease in picture quality caused by unfilled pixel cells.\n2. Discussion of the Related Art\nA liquid crystal display (LCD) device controls an optical transmission ratio of liquid crystal molecules using an electric field to thereby display an image. In general, an LCD device is provided with a liquid crystal panel having liquid crystals between two glass substrates, a matrix of liquid crystal cells, with switching devices for respectively changing signals in the liquid crystal cells, a driving circuit for driving the liquid crystal panel, and a backlight unit for directing light onto the liquid crystal panel.\nRecently there have been efforts to reduce the number of signal lines or the number of circuit components of the liquid crystal panel to develop a thin and light, as well as low cost liquid crystal display device. For example, Korean patent publication No. 10-2003-0039972 discloses a flexible printed circuit and a liquid crystal display device without a gate printed circuit board that can reduce the panel size and the production cost of the liquid crystal display device.\nHowever, because a plurality of data and gate driver integrated circuits, data printed circuit boards, and the like are provided, the flexible printed circuit and the liquid crystal display device even without having a gate printed circuit board have problems in that the reduction of the size and the production cost is limited, and power consumption is high.\nFor the related art LCD device, if the ambient temperature of liquid crystal panel is low (for example, below \u22125\u00b0 C.), the charging property of pixel cell is lowered so that it causes the deterioration of picture quality due to unfilled pixel cells."}
-{"text": "1. Field of the Invention\nThe present invention relates to network elements capable of providing new and enhanced broadband, information-related services to consumers. Specifically, the present invention provides the ability to receive and respond to diverse subscriber requests by transmitting information of virtually any kind. The information transmitted can include movies, video games and other entertainment programs, educational information and programs, scientific and other research-related database information, and the like, and is hereinafter referenced generically as \"information\" or \"titles \".\n2. Description of the Related Art\nConsumer demand for enhanced on-site entertainment and information services is on the rise. Examples of such entertainment and information services include so-called \"on-demand\" video, interactive video games, database research, \"home-shopping\" and the like. Numerous service providers are currently eager to tap this demand, and are therefore expressing interest in schemes for providing such services to consumers. The efforts currently being expended by service providers, as well as the current state of progress in this field of technology are described, e.g., in TIME Magazine, Apr. 12, 1993, at 49ff., and Newsweek, May 31, 1993, at 39ff.\nIn order to be economically practical and viable, a system offering the wide range of information-related services just described should preferably be able to (i) permit access to large volumes of dam (i.e., to a large number of titles), (ii) service a large number of consumers, and (iii) transfer data (i.e., titles) to the end user quickly, preferably in \"real time\". Further, the consensus prefers a system arranged as a centralized facility capable of communicating with remote stations (e.g., homes, offices, public gathering locations, etc.) via existing or next-generation communications media.\nParticular on-going efforts related to the present invention include a system being developed by Broadband Technologies. General elements of the Broadband Technologies system, as understood, are illustrated in FIG. 15. As shown, a subscriber's remote location 151 is connected over a standard telecommunications network 152 to a central facility 153. The central facility 153 awaits service requests from its subscribers.\nWhen a given subscriber wishes to view a selection (e.g., a movie), he or she places a telephone call to the central facility using a standard telephone 154. The call is routed via standard telephone switches 155 (e.g., EO and AT switches) over the telephone voice path of a standard telephone line 156 (e.g., 24-gauge twisted pair copper wire). The telephone call is directed to a Voice Response Unit (VRU) 157 located in the central facility. The VRU forwards the subscriber request to an on-site attendant 158, e.g., through a paper printout. The operator 158 reviews the request, then retrieves the requested selection from an in-house library 159. Finally, the operator 158 locates an available playback unit 170 in a playback battery 171 and loads and plays the selection accordingly.\nThe output signal of the playback unit 170 is supplied to an interface unit 172 coupled between the playback unit 170 and the VRU 157 in the central facility 153. The interface 172 compresses the signal received and forwards the compressed signal via the VRU 157 to the telephone voice path of the telephone line 156. From there, the signal is routed to the subscriber's telephone 154, where a further interface 173 retrieves the signal, decompresses it, and passes it to the subscriber's television set 174.\nAs such, the Broadband Technologies system succeeds in providing an adequate transfer of video data over existing switches and transmission lines. First, the signal received at the remote location 151 corresponds in quality approximately to that of a conventional slow-play VHS mode signal. Furthermore, this system, while not truly \"interactive\", does provide limited system responses to remotely transmitted user requests. However, the system suffers under a range of inefficiencies, limitations on size and flexibility, and limitations or interactiveness, discussed in greater detail below. Finally, the system requires direct linkage to the subscriber's telephone 154, which monopolizes the telephone line and requires electrical connection between the telephone and television.\nAn additional pilot system being developed, known as the AT&T Leesburg Project, overcomes at least some of the problems noted above. General elements of the Leesburg system, as understood, are illustrated in FIG. 16. As apparent from the figure, the Leesburg system is very similar to the Broadband Technologies system, except that the request and response signals are not routed via standard telephone lines. Instead, the request and the response transmission are routed between the remote location 161 and the central facility 163 via a fiber-optic cable system 162.\nIn operation, a subscriber request is initiated via a customized telephone or control unit 164 connected to a special cable converter box 165 installed at the remote location 161. The request is transmitted over the fiber system 162 to a Voice Response Unit (VRU) 166 located at the central facility 163.\nThe operations at the central facility are the same as those described above with respect to the Broadband Technologies system, as indicated by the use of like reference numbers. The output signal of the playback unit 170 is supplied to an interface unit 167 directly coupled to the fiber system 162. The special cable converter box 165 at the remote end retrieves the signal and passes it on to the subscriber's television 168.\nSince the Leesburg system benefits from the larger data transmission capacity of fiber-optics over standard copper wire (approx. 150 Mbits/sec versus approx. 1.5 Mbits/sec), the resulting quality of the data signal is enhanced relative to that of the signal delivered by the Broadband Technologies system. Moreover, this system does not monopolize the subscriber's telephone line. However this system fails to overcome the other problems noted, namely inefficiency, limitations on size and flexibility, and limitations on interactiveness.\nThe problems inherent in these systems, mentioned briefly above, will now be described more particularly.\nThe first class of problems involves inefficiency. First, no system yet proposed is fully automated. The human involvement and interaction required renders such systems costly, fallacy-prone and sluggish in response time. Second, pronounced redundancy between subscribers and playback capacity is required. In other words, the system must dedicate one VCR per concurrent user per average title duration. Further, the system must stock multiple copies of each title in order to cover overlapping requests. Clearly, this redundancy adds tremendously to overall cost and complexity of the system. Moreover, the arrangement is very wasteful, given that certain selections (e.g., first-run movies) are extremely popular for a time and then drop dramatically in popularity.\nThe second class of problems involves limitations on size and flexibility of the overall system. More specifically, given the redundancies discussed above, any system designed in this manner is inherently limited in the number of subscribers that can be accommodated and in the number of selections (titles) that can be offered.\nThe third class of problems relates to interactiveness. In each of the systems described, the user merely requests and views the title selected. As designed, the systems are incapable of providing more versatile or elaborate options that would require greater interaction between the user and the system. Examples of possible options available from a more powerful system include obtaining information to assist the user in selecting a title, selecting and reviewing previews of titles, adjusting playback of titles (e.g., pause, fast-forward, etc.), selecting between original language and dubbed versions of titles, etc. Further, truly interactive services are foreclosed from the described systems. Such interactive services of interest to consumers include on-line database searching, networked remote-station video-game competitions, personalized home shopping, and the like."}
-{"text": "It is well known in the vehicle industry that certain rotary accessory devices are often used in connection with vehicle motors. Two such accessory devices are alternators and starter motors. Alternators are used in connection with an engine and are typically belt driven by the engine. Alternators have internal components, which when rotated supply electrical power to a vehicle battery. Alternators are typically removable but rigidly mounted via a bracket to the engine block or the chassis of the vehicle. In many cases, where a standard type of alternating mounting arrangement is used, the alternator has \u201cears\u201d with holes that are mounted onto a post or bolt attached to the vehicle. This permits pivoting of the alternator so that the alternator can be pivoted around the post against the belt tension in order to install and remove belts, and provide a suitable tension when the belt is installed.\nStarter motors are electrical motors, which are typically removable but rigidly mounted to an engine or transmission casing. The starter motor has an electrically driven pinion gear extending from the starter motor that engages a component (typically gears on the flywheel of the engine) in order to be able to rotate the crankshaft of the engine to start it. There is a wide range of attachment mechanisms for attaching the described starter motor.\nConventional alternator and starter tester may test alternators and/or starter motors using a direct current (DC) power source. For example, the DC power source may be used to power the starter motor and bias the alternator during a testing process. The DC power source may determine the capability of an alternator and starter tester. Oftentimes, the alternator and starter tester requires high power in order to achieve desired accuracy when testing the starter motor or the alternator. The high power required leads to expensive and bulky circuit components which in turn lead to higher cost, size and weight for the alternator and starter tester. Additionally, the high power demand required by the prior art systems resulted in greater power consumption and operating costs. Thus, it would be desirable to have an alternator or starter tester with high accuracy with decreased cost, size and weight."}
-{"text": "The invention relates to processing multiview video.\nMultiview Video Coding (MVC) relates to compression of video sequences (e.g., a sequence of images or \u201cpictures\u201d) that are typically acquired by respective cameras. The video sequences or \u201cviews\u201d can be encoded according to a standard such as MPEG. A picture in a video sequence can represent a full video frame or a field of a video frame. A slice is an independently coded portion of a picture that includes some or all of the macroblocks in the picture, and a macroblock includes blocks of picture elements (or \u201cpixels\u201d).\nThe video sequences can be encoded as a multiview video sequence according to the H.264/AVC codec technology, and many developers are conducting research into amendment of standards to accommodate multiview video sequences.\nThree profiles for supporting specific functions are prescribed in the current H.264 standard. The term \u201cprofile\u201d indicates the standardization of technical components for use in the video encoding/decoding algorithms. In other words, the profile is the set of technical components prescribed for decoding a bitstream of a compressed sequence, and may be considered to be a sub-standard. The above-mentioned three profiles are a baseline profile, a main profile, and an extended profile. A variety of functions for the encoder and the decoder have been defined in the H.264 standard, such that the encoder and the decoder can be compatible with the baseline profile, the main profile, and the extended profile respectively.\nThe bitstream for the H.264/AVC standard is structured according to a Video Coding Layer (VCL) for processing the moving-image coding (i.e., the sequence coding), and a Network Abstraction Layer (NAL) associated with a subsystem capable of transmitting/storing encoded information. The output data of the encoding process is VCL data, and is mapped into NAL units before it is transmitted or stored. Each NAL unit includes a Raw Byte Sequence Payload (RBSP) corresponding to either compressed video data or header information.\nThe NAL unit includes a NAL header and a RBSP. The NAL header includes flag information (e.g., nal_ref_idc) and identification (ID) information (e.g., nal_unit_type). The flag information \u201cnal_ref_idc\u201d indicates the presence or absence of a slice used as a reference picture of the NAL unit. The ID information \u201cnal_unit_type\u201d indicates the type of the NAL unit. The RBSP stores compressed original data. An RBSP trailing bit can be added to the last part of the RBSP, such that the length of the RBSP can be represented by a multiple of 8 bits.\nThere are a variety of the NAL units, for example, an Instantaneous Decoding Refresh (IDR) picture, a Sequence Parameter Set (SPS), a Picture Parameter Set (PPS), and Supplemental Enhancement Information (SEI), etc.\nThe standard has generally defined a target product using various profiles and levels, such that the target product can be implemented with appropriate costs. The decoder satisfies a predetermined constraint at a corresponding profile and level.\nThe profile and the level are able to indicate a function or parameter of the decoder, such that they indicate which compressed images can be handled by the decoder. Specific information indicating which one of multiple profiles corresponds to the bitstream can be identified by profile ID information. The profile ID information \u201cprofile_idc\u201d provides a flag for identifying a profile associated with the bitstream. The H.264/AVC standard includes three profile identifiers (IDs). If the profile ID information \u201cprofile_idc\u201d is set to \u201c66\u201d, the bitstream is based on the baseline profile. If the profile ID information \u201cprofile_idc\u201d is set to \u201c77\u201d, the bitstream is based on the main profile. If the profile ID information \u201cprofile_idc\u201d is set to \u201c88\u201d, the bitstream is based on the extended profile. The above-mentioned \u201cprofile_idc\u201d information may be contained in the SPS (Sequence Parameter Set), for example."}
-{"text": "1. Field of the Invention\nThe present invention relates to a luminescent storage screen of the type suitable for latently storing an x-ray image in an x-ray diagnostics installation, and an apparatus for reading-out the stored image from the screen.\n2. Description of the Prior Art\nAn x-ray diagnostics installation is described in German OS No. 35 29 306 which includes a luminescent storage screen having a photo stimulable phosphor which is irradiated by x-rays attenuated by an examination subject. The phosphor serves as a radiation-sensitive transducer in which holes (in the electronic sense) are generated, and are stored in potential traps or wells. The entire surface of this storage screen is scanned pixel-by-pixel during read-out by a different radiation source, for example a laser, so that the holes stored in the traps are excited, and then fall back into a stable energy level. The energy difference is emitted in the form of light quanta. The x-ray image latently stored in the screen can thus be read-out of the screen at a later time.\nThe storage screen in this known installation consists of two phosphors separated by a light-impermeable layer. These phosphor layers are scanned surface-wide by a beam deflected by a mirror, for example, in the vertical direction, with the beam being divided by a beam splitter so that both phosphor layers are simultaneously scanned. All picture elements in a line of the storage screen are successively scanned by this deflection. The entire screen surface is scanned by moving the storage screen perpendicularly relative to the line scanning, advancing the screen line-by-line. The light emitted by the stimulable phosphor layers as a result of the interaction with the read-out beam is acquired by two optical collectors, and is supplied to the light-sensitive input faces of two separate light detectors. The respective output signals of the detectors are supplied to a video chain for reproduction on a display, with the output signals of the two detectors being combined.\nInstead of dividing the read-out beam, however, it is also possible to use two separate radiation generators, each undertaking a separate scanning of one of the phosphor layers. In this case, however, it is necessary to synchronize the deflection of both scan beams, so that the picture elements of the two phosphor layers which correspond to each other are simultaneously read-out. Synchronization need not be undertaken at the time of read-out; it can alternatively be undertaken electronically in a processing circuit following the light detectors.\nThe two storage layers in combination function as a single storage layer having an increased thickness. It is known that the efficiency of a phosphor layer increases with increasing thickness, since more x-ray radiation is absorbed in a thicker layer. Image sharpness, however, decreases with increasing layer thickness, so that fine details of the image can no longer be adequately portrayed. By using two separate phosphor layers, with a separate read-out of each layer, the efficiency of a layer having a thickness corresponding substantially to the combined thickness of the two layers is achieved, while a better image sharpness is obtained than would be possible using a single layer of the same combined thickness. A disadvantage is the complexity of the read-out system which is required for such a screen, which not only requires two separate read-out beams, but also structure for guiding and synchronizing the two beams so that the pixels separately read-out from each phosphor layer can be correlated to form a single image. Two synchronized scan systems, or at least two synchronized scan beams, are required.\nAnother x-ray diagnostics installation is described in OS No. 29 51 501 which includes a luminescent storage screen having a phosphor layer scanned by a read-out beam, with detectors being disposed on both sides of the storage screen. These detectors respectively collect the light emitted from the phosphor layer in an upward direction and the light emitted in a downward direction. A higher light yield and thus an improved bright-dark contrast are thereby achieved. Because only a single phosphor layer is used, however, the competing problems of efficiency of x-ray radiation absorption and image sharpness are still present."}
-{"text": "This invention relates to method and apparatus for decoding PCM digital input words, and more particularly to an integrated PCM decoder which is substantially insensitive to parasitic and stray capacitance effects.\nPCM encoders and/or decoders are described in the articles \"A Two-Chip PCM Voice Codec With Filters\" by Y. A. Haque, et al., IEEE Journal of Solid State Circuits, Vol. SC-14, No. 6, pages 961-969, December 1979; \"A PCM Voice Codec With On-Chip Filters\" by J. T. Caves, et al., IEEE Journal of Solid State Circuits, Vol. SC-14, No. 1, pages 65-73, February 1979; \"A Segmented Mu-255 Law PCM Encoder Utilizing NMOS Technology\" by Paul R. Gray, et al., IEEE Journal of Solid State Circuits, Vol. SC-11, No. 6, pages 740-747, December 1976; and \"A Unified Formulation of Segment Companding Laws and Synthesis of Codecs and Digital Compandors\" by H. Kaneko, The Bell System Technical Journal (BSTJ), September 1970, pages 1555-1588, which are incorporated herein by reference. Mu-law decoders operating on 8-bit PCM code words have employed arrays of binary weighted capacitors, the total array capacitance typically being 255 times the capacitance of the smallest capacitance of the array. The capacitance of the two largest binary weighted capacitors of the array are therefore 128 and 64 times that of the smallest capacitor thereof. In integrated circuit structures it is desirable to minimize the surface area of a chip that is required for individual circuit elements. It is readily seen that the chip area dedicated for integrated capacitors in such an integrated decoder may be reduced by approximately 50% by deleting only one of the binary weighted capacitors. An integrated decoder requiring a capacitance of only 32 times that of the smallest capacitor of an array is disclosed in copending patent application (infra).\nAn object of this invention is the provision of an improved PCM decoder requiring only three equal valued capacitors."}
-{"text": "Ocular surgery, whether reconstructive, cosmetic, palliative, or otherwise, is highly patient specific. Even though most surgery patients have the same basic ocular architecture, every eye has its own set of specific features and dimensions that in certain cases may be significantly different from those of expected norms. As a result, surgeons must rely upon their individual experience and skills to adapt whatever surgical techniques they are practicing to the individual requirements as determined by each patient's unique ocular structural features and dimensions.\nTo date, this individualized surgical adaptation has been accomplished essentially through freehand and best guess techniques based upon a pre-surgery examination and evaluation of the individual patient's ocular region. This examination may include preliminary measurements as well as the surgeon making reference markings directly on the patient's ocular tissues with a pen or other form of dye or ink marking. Then, after the patient has been prepared and placed in position for surgery, typically in a supine or prone position as opposed to the often vertical positioning of the patient during the pre-surgery examinations, the surgeon adapts the placement and configuration of the initial surgical incisions to the actual physical dimensions and circumstances found in the patient as the surgical procedure progresses. As a result, many initial measurements or reference markings on the patient's ocular tissues are at best a general guide as to where to begin the procedure and have limited accuracy and influence on subsequent aspects of the procedure or on the overall outcome of the surgery.\nFurther complicating matters, ocular tissues are not conducive to pre-surgery reference markings or measurements. This is particularly true because most ocular tissues have wet surfaces and internal physical structures that cannot be accessed for direct measurement or marking prior to surgery.\nAdditionally, pre-surgical rinsing, sterilization, or drug administration to the ocular tissues may dissolve, alter or even remove reference markings prior to the initiation of surgery. Similarly, subsequent wiping and contact with fluids, including the patient's body fluids, during the surgical procedure may remove or distort any remaining reference markings. As a result, even the most accurate surgical reference markings may lose any practical effectiveness beyond the initial stages of the surgical procedure.\nAs such, there is a continuing need for effective surgical reference indicia properly aligned with a particular axis, particularly when proper alignment of pre-surgical data is pivotal to satisfactory patient outcome. For instance, proper alignment of pre-surgical data, particularly the true vertical axis of an eye, with the ocular surgery is highly advantageous with intraocular lens (IOL) implantation and orientation of the IOL within the posterior chamber of the eye.\nThe implantation of an IOL is a highly sophisticated surgical procedure that, in the past, has been performed based on partially accurate or even inaccurate visual measurements and estimated alignment with the vertical or horizontal axis of a patient's eye. Past procedures have commonly relied on the measurement of an eye's vertical or horizontal axis prior to surgery and the subsequent inaccurate translation of that measurement to the IOL implantation procedure where the positioning of the measured axis of the eye may have changed. As a result, it is not uncommon for the IOL placement to be improperly aligned with the true vertical axis of the eye, resulting in such side effects as poor visual acuity, double vision, and halos under low ambient light conditions.\nAccordingly, in spite of the ongoing development and the growing sophistication of contemporary ocular surgery, there is a continuing need for the provision of effective surgical reference indicia including information about or which track a patient's natural vertical axis or other important axis of orientation."}
-{"text": "The present invention relates to evaporative emission control systems for vehicles and in particular to a purge valve that is adapted to be controlled by the engine management control system for regulating the supply of fuel vapors to the engine intake from the fuel tank vapor recovery system.\nIn order to meet current emission requirements, present day vehicles contain evaporative emission control systems which reduce the quantity of gasoline vapors emanating from the fuel tank of the vehicle. Generally, these systems include a charcoal canister which traps the vapors from the fuel tank, and a purge system which draws the vapors out of the canister and feeds them into the intake system of the engine when the engine is running. The fuel vapors are drawn into the engine intake manifold along with atmospheric air drawn through the canister.\nThe capability of the canister to trap vapors from the fuel tank is greatly dependent upon how thoroughly the vapors are purged from the canister when the vehicle was last operated. Accordingly, it is desirable to purge the canister as much as possible while the engine is running. However, the amount of vapor that can be drawn into the engine at any time is limited by the total airflow into the engine and the accuracy with which the purge flow can be controlled. At high speeds or under high engine loads, high purge flow rates can be easily handled. Under such conditions, however, the manifold vacuum is low which tends to limit the amount of fuel vapors and air which can be drawn from the canister into the engine intake manifold. In addition, when the engine is at idle, the airflow into the engine is low. Therefore, purgining at idle must be precisely controlled to prevent a rough idle. Moreover, due to the varying ratio of air to fuel vapors in the purge system, purging during idle can significantly impact the resulting air/flow ratio of the fuel mixture supplied to the engine. Consequently, purging at idle can easily result in a too rich or too lean fuel mixture causing excessive tailpipe emissions unless purging at idle is limited to low flow rates. Current emissions systems, therefore, do not generally purge the canister at idle to any substantial degree.\nHowever, impending tighter emissions requirements and changes to the EPA testing procedures will require larger capacity canisters and therefore higher capacity purge systems. Moreover, the prospect of on-board refueling vapor recovery systems will only add to these system requirements. Accordingly, it is becoming imperative that such systems not only purge at idle, but that maximum flow rates be increased as well. This, of course, presents conflicting requirements for purge systems. Specifically, in order to purge at idle, the purge flow rate must be fairly low and accurately controlled by the engine control computer which monitors the resulting oxygen content of the exhaust gases from the engine. When a canister is saturated with fuel, and vapor is initially purged, the purge flow is very high in fuel vapor. After most of the fuel vapors are drawn out of the charcoal, the purge flow is almost pure air. Therefore, the purge control valve must be capable of allowing the engine control computer to precesely control small flow rates at idle while correcting the idle fuel-air ratio so that tailpipe emissions are not adversely affected. This type of precise flow control is best accomplished using a relatively small valve.\nOn the other hand, it is desirable to purge at very high flow rates when the engine is operating under high speed or heavy load conditions when it can efficiently consume significant quantities of fuel vapor and air with a minimum effect on fuel air ratios. In order to achieve large flow rates, it is necessary for the purge valve to provide a relatively large flow passage. This requirement, of course, is in direct conflict with the requirement for precise low flow rate control. Specifically, it is believed to be impractical to provide a valve large enough to satisfy the high flow requirements which at the same time is capable of precisely modulating the opening of the valve to meet the low flow requirements.\nAccordingly, it is the primary object of the present invention to provide a two-stage purge control valve that is capable of providing both precise control at low flow rates and high flow capacity at low manifold vacuum pressures. In general, this is accomplished by providing a single assembly having two valves which control separate parallel flow paths. Low flow control is achieved with a small solenoid valve adapted to be driven by a pulse width modulated (PWM) signal from the engine control computer. High flow capacity is provided by a vacuum-controlled valve which opens at low manifold vacuum pressures. Because purge flow comprises a relatively small percentage of total air flow into the engine under the conditions when the high flow stage is open, precise control of the high flow capacity valve by the engine control computer is not required.\nAccordingly, the purge valve according to the present invention allows the full range from 10% to 90% duty cycle control to be used to control low flow rates and opens the high flow valve only when the purge flow comprises a small portion of the total engine intake air flow. Moreover, the high flow valve is adapted to open gradually as engine manifold vacuum pressure decreases, thereby proportioning the purge flow to the total engine intake air flow. In addition, the engine control computer can still adjust the high purge flow rate to a degree by controlling the parallel flow through the PWM solenoid valve.\nIn the preferred embodiment of the present invention, the response and flow capacity of both the low and high flow control valves can be calibrated to meet the requirements of a particular engine family or purge system."}
-{"text": "This application claims priority to Application No. 10026240.6 which was filed in the German language on May 26, 2000.\nThe invention relates to a method for optical polarization control, and in particular, to endless polarization control of an optical signal.\nFor optical polarization control, polarization transformers are particularly suitable, preferably those of an electro-optical principle of operation. These transformers typically include one or more polarization actuators, the modes of which can be varied endlessly, i.e. without interruption, on a great circle of the Poincarxc3xa9 sphere. Examples of this are known from IEEE J. Quantum Electronics 18(1982)4, pp. 767-711, from IEEE J. Lightwave Techn. 6(1988)7, pp. 1199-1207, from IEEE J. Lightwave Techn. 8 (1990) 3, pp. 438-458, from Electron. Lett. 27(1991)4, pp. 377-379, and from the minutes of the European Conference on Optical Communications 1993, Montreux, Switzerland, pp. 401-404, contribution WeP9.3.\nIn practice, non-ideal behavior of polarization actuators can impede the capability for endless, uninterrupted polarization tracking which is always the aim. In IEEE J. Lightwave Techn. 8(1990)3, pp. 438-458, it was specified how this problem can be solved purely quantitatively, i.e. by adding another electro-optical polarization actuator. More elaborate tests have shown, however, that only one additional polarization actuator is hardly sufficient in practice.\nIn IEEE J. Lightwave Techn. 6(1988)7, pp. 1199-1207, a solution to solving this problem is to use a calibration table in a manner which provides the maximum control speed per control step. This can be disadvantageous, however, in that the constructional length of the polarization transformer normally required is distinctly increased by adding a further polarization actuator to an original electro-optical polarization actuator which sufficient in the ideal case, and is even double in this case. As an alternative, the control voltages could be distinctly increased while keeping the total construction length unchanged. Both methods contradict the usual requirements for low supply voltages and short constructional length for the purpose of avoiding DC drift and optical insertion losses. Using this technique, the further polarization actuator is driven as a function of the parameters of the original one. In the nomenclature used there, a phase angle difference, occurring for d3=0, d4=0, of d2xe2x80x2=atan2(f4(d2), f3(d2)) (atan2=arcfunction, known for example from the programming languages Pascal, C, Matlab) of the further actuator which occurs in the case of input polarization P1, which is circular here, between this and a proportion of output polarization P2 (which is orthogonal thereto and is thus oppositely circular here), is also a function of a corresponding phase angle difference d2 of the original one which is accounted for by a calibration table. Since figure 15 there is a twisted curve, the problem arises that atan2(f4(d2),f3(d2)) and d2 can differ, for example by an odd-number multiple of xcfx80. In these cases, the total delay of the polarization transformer formed by the two polarization actuators is obtained by subtracting the delay d1 of the original one and the delay sqrt (f3(d2){circumflex over ( )}2+f4(d2){circumflex over ( )}2) of the further polarization actuator. That is, d1xe2x88x92sqrt(f3(d2){circumflex over ( )}2+f4(d2){circumflex over ( )}2) for d1 greater than sqrt(f3(d2){circumflex over ( )}2+f4(d2){circumflex over ( )}2). This destructive interplay is the reason why a large constructional length and/or high control voltages are needed in accordance with the prior art. If the method specified is properly performed, atan2(f4(d2),f3(d2)) will change by 4xcfx80 or d2 by 2xcfx80 so that the amount |atan2(f4(d2),f3(d2))xe2x88x92d2 of the difference atan2(f4(d2),f3(d2))xe2x88x92d2 can grow without limits if d2 grows arbitrarily, e.g. by many times 2xcfx80. In particular, the discussed destructive interplay occurs there, for example with d2=0.7*xcfx80, d2=1.3*xcfx80 and d2=1.8*xcfx80.\nSince this phase angle difference (d2 or atan2(f4(d2),f3(d2))) in each case represents an angle coordinate from an eigenmode of the polarization actuator on a great circle on the Poincarxc3xa9 sphere, the words angle coordinate will be used synonymously with this phase angle difference in the text which follows.\nIn IEEE J. Lightwave Techn. 6(1988)7, pp. 1199-1207, an electro-optical polarization transformer for transforming a particular one into any arbitrary polarization state or conversely has been described; such polarization transformers could be abbreviated by GSBA (Generalized Soleil Babinet Analogue). Polarization transformers for transforming any arbitrary polarization state into any arbitrary one will be abbreviated by ER (Elliptical Retarder) in the text which follows. ERs with calibration table(s) for compensating for non-ideal component behavior have not previously been known. This is especially true of ERs that necessitate a short constructional length since, compared with GSBAs, ERs already need twice the constructional length in the ideal case.\nIn Proc. Fourth European Conference on Integrated Optics ECIO 87, Glasgow, Scotland, pp. 115-118, a GSBA is specified with many sections and fixed angle coordinate differences (For example: xcex1*xcfx80/2 between adjacent electrodes, where xcex1=xcex9DEV/xcex9BEAT), which are only dependent on the wavelength but not on the desired polarization transformation, of different sections. Correspondingly, this GSBA also does not have a calibration table for compensating for non-ideal behavior.\nIn one embodiment of the invention, there is a method for endless polarization control of an optical signal. The method includes, for example, converting a polarization of the optical signal using polarization acutuators, with arbitrary, freely and interruption free selectable individual phase angle differences as individual angle coordinates, at least partially into a respective orthogonal, which together form a polarization transformer with a total delay and a total phase angle difference as a total angle coordinate between the polarization and a respective orthogonal, in which an individual angle coordinate is varied as a function of another angle coordinate, wherein the amount (|dxcex61|, |dxcex62|, |dxcex6|, |dxcex6*(2*ixe2x88x92nxe2x88x921)/(nxe2x88x921)|) of the variation (dxcex61, dxcex62, dxcex6, dxcex6*(2*ixe2x88x92nxe2x88x921)/(nxe2x88x921)) of an individual angle coordinate is limited for sequences of desired polarization transformations.\nIn another aspect of the invention, individual delays support each other in their effects ((\"psgr\"1xe2x80x2 greater than \"psgr\"1, \"psgr\"xe2x80x2 greater than \"psgr\"2), (\"psgr\" greater than \"psgr\"1, \"psgr\" greater than \"psgr\"2)) by the limitation of the amount (|dxcex61|, |dxcex62|, |dxcex6|, |dxcex6*(2*ixe2x88x92nxe2x88x921)/(nxe2x88x921)|) of the variation (dxcex61, dxcex62, dxcex6, dxcex6*(2*ixe2x88x92nxe2x88x921)/(nxe2x88x921)).\nIn another aspect of the invention, the another angle coordinate is the actual or attempted total angle coordinate or an attempted individual angle coordinate.\nIn yet another aspect of the invention, a number of individual angle coordinates are varied as a function of the another angle coordinate.\nIn another aspect of the invention, at least in the case of equal attempted individual angle coordinates, one of the variations of individual angle coordinates is equal to the negative of another one of the variations or equal to zero.\nIn another aspect of the invention, at least one of the individual delays is varied as the function of another angle coordinate.\nIn still another aspect of the invention, one of the variations of the individual delays is equal to another one of the variations or equal to zero, at least in the case of equal attempted individual delays.\nIn another aspect of the invention, the method includes applying, to at least two additional polarization actuators, which are located before or after the polarization actuators in the beam path to produce another polarization transformer and form an elliptical retarder.\nIn another aspect of the invention, one of the polarization actuators operates as an electrooptical Soleil Babinet compensator which can convert circular polarizations at least partially into one another.\nIn yet another aspect of the invention, at least two additional polarization actuators producing an additional polarization transformer form an elliptical retarder with the polarization transformer.\nIn another aspect of the invention, at least one of the polarization actuators operates as an electrooptical Soleil Babinet compensator, the polarizations which are at least partially converted being circular polarizations.\nIn another aspect of the invention, at least one of the polarization actuators operates as an electrooptical Soleil Babinet analog, the polarizations which are at least partially converted being TE and TM polarizations."}
-{"text": "Balloon catheters are increasingly being employed to conduct neurological procedures in patients. However, the design parameters for balloon catheters intended for use in neurological procedures are significantly different than the design parameters for balloon catheters used in non-neurological procedures such as cardiological procedures. For example, the width of the circulatory system within the neuroanatomy is significantly smaller and more tortuous than the circulatory system in other parts of the body. In order to access the smaller and more tortuous regions of the neuroanatomy, it is necessary to minimize the outer diameter of the balloon catheter while simultaneously maintaining the pushability and trackability of the catheter.\nIn order to minimize the outer diameter, current balloon catheters intended for neurological procedures employ a non-reinforced, single lumen, over-the-wire design. Accordingly, these balloon catheters are prone to several problems. First, the non-reinforced lumen is susceptible to ovalizing and/or kinking which, in turn, hinders advancement of the catheter over the guidewire, as well as deflation of the balloon. Second, the single lumen is in communication with the arterial blood flow. As the guidewire and balloon catheter are manipulated through the circulatory system, blood is withdrawn into the single lumen of the balloon catheter. Blood may thereby enter the balloon during inflation and cause (1) poor imaging of the balloon, for example, poor fluoroscopic imaging; (2) poor passage of the balloon through the circulatory system due to the premature inflation of the balloon; and (3) poor deflation of the balloon due to blood coagulation in the balloon inflation/deflation port. An additional disadvantage of single lumen balloon catheters is that the interference fit of the guidewire and inflation seal of the balloon may result in removal or peeling of the hydrophilic coating of the guidewire.\nIn order to minimize the outer diameter, current balloon catheters intended for neurological procedures are also typically designed to work with only a narrow gauge guidewire that is supplied by a manufacturer along with the balloon catheter. The current balloon catheters employ guidewires having diameters in the range of 0.010 to 0.012 inches. These relatively narrow guidewires are soft and, therefore, are very difficult to maneuver through the small, tortuous neuroanatomy.\nWhat is needed in the field is a balloon catheter that is operable to use with larger gauge guidewires; resists ovalizing and kinking of the inflation and guidewire lumen(s); and deploys with improved pushability and trackability."}
-{"text": "The present invention relates to a router bit. A router bit used for chamfering and finishing the butt end and side of a board is available in two types, i.e. a straight tooth type and a helical tooth type. Usually, a router bit of the helical tooth type is provided with an even number of cutting edges having helix angles in the same hand or direction with respect to the axis of rotation of the cutting edges. The conventional router bit of this type is liable to form burrs at the side of board when a board material containing fiber reinforced plastics is chamfered or finished. This is true particularly in chamfering or finishing FRP board or FRP-clad board in which fluffy fibers remain on the finished surface because the fibers are much stronger than the matrix resin. Such rough finish brings about delamination, corrosion, a decrease of strength, and breakage. This is quite different from problems involved in chamfering or finishing metal materials.\nThe object of the present invention is to eliminate the drawbacks of conventional technology as mentioned above, and is the result of observation of mutual actions of cutting edges and the material being cut."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to hard-coated antiglare films, and polarizing plates and image displays including the same.\n2. Description of the Related Art\nWith technical improvements in recent years, liquid crystal displays (LCDs), plasma display panels (PDPs), electroluminescence displays (ELDs), etc. have been developed in addition to conventional cathode ray tube (CRT) displays as image displays and have been used in practical applications. Particularly, with technical innovations of LCDs with respect to, for example, wide viewing angles, high resolution, high response, and good color reproduction, applications of LCDs are spreading from laptop personal computers and monitors to television sets. In an LCD, generally, a liquid crystal panel is used in which polarizing plates are disposed on both sides of a liquid crystal cell, respectively. Generally, in order to prevent scratches on the polarizing plate, the surface of the liquid crystal panel is subjected to a hard coating treatment. For the hard coating treatment, a hard-coated film is often used. The hard-coated film is subjected to an antiglare treatment for preventing a decrease in contrast due to reflection of outside light such as fluorescent light and sunlight as well as reflection of an image at the surface of the liquid crystal panel. Particularly with an increase in size of the screens of image displays, the number of image displays equipped with hard-coated antiglare films is increasing.\nRecently, in order to improve image quality, the number of high definition image displays with small pixel sizes is increasing. In such high definition image displays, when a conventional hard-coated antiglare film is disposed, variations in brightness in pixels are emphasized further more to cause a visible failure (a failure due to glare) to result in considerably deteriorated image quality. Conventionally, in order to produce an antiglare laminate that allows high definition to be provided, a method has been employed in which the haze value of the antiglare layer is increased to eliminate glare. However, this method has a problem in that the contrast is decreased considerably because light is scattered intensively at the panel surface. Furthermore, in the case where the surface unevenness is increased to improve antiglare properties, there is a problem in so-called white blur in an oblique direction, that is, white blur is observed due to excessively intensive scattering of reflected light when the panel surface is viewed from the oblique direction. In the antiglare treatment, a method of producing an uneven shape of the film surface by adding, for example, inorganic or organic particles is used. Generally, although an improvement in antiglare properties and an improvement in contrast or white blur are considered to have a contradictory relationship, various proposals have been made to obtain both these properties. For example, it is being studied that aggregates with three-dimensional structures formed of the aforementioned particles are allowed to be present in an antiglare layer (see, for instance, JP 2005-316413 A). This, however, may cause scatter at the aggregates or may cause a fine pattern to appear in a hard-coated film. Furthermore, although a method that is effective in improving a part of properties has been proposed (see, for instance, JP 2003-4903 A and JP 4001320 B), a method that is effective in solving all the three problems described above has not been found."}
-{"text": "1. Field of the Invention\nThe present invention is directed to a flow transducer muff for positioning a substantially right-angled parallelepiped-shaped flow transducer in the expiration channel of a ventilator, of the type having first and second muff halves which are connectable on opposite sides of the flow transducer.\n2. Description of the Prior Art\nExpired gas flow is measured in a flow transducer at the expiration side of a ventilator. The flow transducer is maintained in position in the expiration channel by means of a flow transducer muff, which forms a gas-tight connection to the flow transducer. Such a flow transducer muff is described, for example, in the Operating Manual (page 10: 2-3, August 1988) for the Servo Ventilator 900C, sold by Siemens-Elema AB. In general, known flow transducer muffs consist of two halves which are connected on opposite sides of the flow transducer.\nThe two muff halves in known components are identical in shape, and there is thus a risk that the flow transducer could be installed, after cleaning, in a reversed position in the ventilator."}
-{"text": "The present invention relates in general to heat dissipation in semiconductors and, more particularly, to direct contact liquid cooling in electronic circuit module packages.\nIt is known in the art that combinations of various type semiconductors and/or hybrid circuits are required to achieve large scale electronic systems. It is desired that the packaging of these system functions be miniaturized in modular form to optimize cost, size, and weight advantages for the end user. The electronic systems often consume an appreciable amount of power and generate heat. As the level of modular system integration increases, the size of the module and the power consumption of the module correspondingly increases.\nThe electronic systems are often used at various pressure altitudes including submerged applications and high altitude applications. In some applications, conventional package cooling techniques such as convection and radiation may be inadequate to assure sufficient cooling and reliable electronic system performance. Likewise, package pressure constraints that are related to a wide range of pressure altitude variations have not been adequately addressed with present module technology to provide a universal pressure solution.\nHence, a need exists for an electronic circuit module package that is capable of reliable operation at various atmospheric and sub-surface altitudes and maintains good heat dissipation characteristics throughout its range of operation."}
-{"text": "1. Field of the Invention\nThe body relies upon a complex immune response system to distinguish self from non-self. The proper functioning of the immune system is vital for the long term health of the body.\nDeficient immune response can lead to the body's inability to protect itself from non-self matter. Excessive immune response can lead to the body's over reaction to what would otherwise be innocuous matter.\nAt times, the body's immune system must be controlled in order to either augment a deficient response or suppress an excessive response. For example, when organs such as kidney, heart, heart-lung, bone marrow, and liver are transplanted in humans, the body will sometimes reject the transplanted tissue by a process referred to as allograft rejection.\nIn treating allograft rejection, the immune system is frequently suppressed in a controlled manner through drug therapy. Immunosuppressant drugs are carefully administered to transplant recipients in order to help prevent allograft rejection of non-self tissue.\nOne drug which finds use as an immunosuppressant in the United States and other countries is cyclosporin A (CsA). CsA is a cyclic undecapeptide of the general structure: ##STR1##\nwherein all of the .alpha.-amino acid residues that form the cyclic backbone of cyclosporin A are of the L-configuration except .alpha.-amino acid 8 which is of the D-configuration. Amino acid residue 1 is derived from an unusual 9 carbon amino acid [2S,3R,4R,6E]-3-hydroxy-4-methyl-2-methylamino-6-octenoic acid. Amino acid residues 1, 3, 4, 6, 9, 10 and 11 are N-methylated on the amide nitrogen atoms of the cyclic backbone of cyclosporin A. Cyclosporin A is described in U.S. Pat. Nos. 4,117,118 (1978) and 4,396,542 (1983).\nCsA may have other useful properties such as antibiotic, anti-arthritic and anti-inflammatory activities and may find use in the treatment of other conditions such as diabetes, malaria and autoimmune diseases.\nA large number of CsA metabolites that retain the undecapeptide ring have been identified and reported (see Maurer, G.; Loosli, H. R.; Schreier, E.; Keller, B. Drug Metabolism and Disposition 1984, 12(1), 120-126, the structures, nomenclature and analytical data of the metabolites are incorporated herein by reference). It is not known what role, if any, these metabolites play in either the desired immunosuppressant activity or any unwanted adverse reactions when CsA is used for the prevention of allograft rejection.\nEven though CsA is a highly effective immunosuppressant drug, its use must be carefully managed because the effective dose range is narrow and excessive dosage can result in serious side effects. Renal dysfunction, hypertension, cardiovascular cramps, hirsutism, acne, tremor, convulsions, headache, gum hyperplasia, diarrhea, nausea, vomiting, hepatotoxicity, abdominal discomfort, paresthesia, flushing, leukopenia, lymphoma, sinusitis and gynecomastia have been observed in kidney, heart or liver transplant patients undergoing CsA treatment. Too little CsA can lead to graft rejection.\nManagement of CsA dosage involves careful control of the level of the drug present in the patient. Because the distribution and metabolism of CsA varies greatly between patients, and because of the wide range and severity of adverse reactions, accurate monitoring of drug level is considered essential.\nLaboratory methods for detection of cyclosporin have been developed. These techniques typically involve high performance liquid chromatography (HPLC), radioimmunoassay (RIA) or florescence polarization immunoassay (FPIA). See for example Wolf, B. A.; et al. Clinical Chem. 1989, 35(1), 120-124; Vernillet, L.; et al. Clinical Chem. 1989, 35(4), 608-611; Ball, P. E.; et al. Clinical Chem. 1988, 34(2), 257-260; Schran, H. F.; et al. Clinical Chem. 1987, 33(12), 2225-2229; Sanghvi, A.; et al. Clinical Chem. 1988, 34(9), 1904-1906; McBride, J. H.; et al. Clinical Chem. 1989, 35(8), 1726-1730; Quesniaux, V.; et al. Clinical Chem. 1987, 33(1), 32-37.\nEach of these techniques has certain limitations with regard to safety and complexity of the procedure and level of specificity for cyclosporins of interest. For example, HPLC requires long sample preparation and/or run times using high cost labor-intensive procedures; RIA presents the well-known hazards of handling radioactive materials; and FPIA, when based on non-specific mono- or polyclonal-antibodies, often fails to distinguish between CsA and its metabolites.\nA simple analytical method specific to selected cyclosporins is needed for use in cyclosporin treatment management.\nImmunoassay is a technique suited to the requirements of a simple method for measuring the amount of cyclosporin in a sample suspected of containing cyclosporin. However, most available antibodies capable of recognizing cyclosporins of interest also recognize and cross-react with closely related compounds such as cyclosporin metabolites. Because of this cross-reactivity, immunoassays dependent on these antibodies are less specific to cyclosporins of interest than might be desired. Cyclosporin immunoassays dependent on antibodies which cross-react with compounds other than cyclosporins of interest can be improved if the cross-reactive compounds are substantially inactivated toward recognition by the antibodies which recognize the cyclosporins of interest.\nThe methods and compositions of the present invention relate to a simple, specific immunoassay method for cyclosporins. The methods and compositions of the present invention also related to a method for inactivating cross-reactive compounds in immunoassay methods for cyclosporin.\n2. Brief Description of the Related Art\nWorld Patent Application No. 8,602,080 (1986), describes monoclonal antibodies selective for certain cyclosporins. These antibodies are prepared in response to cyclosporins bound to antigenic carriers through a modified amino acid residue such as a (D)-Lysine at position 8 or an (L)-threonine at position 2.\nEuropean patent application No. 283,801 (1988), describes a fluorescence immunoassay method and cyclosporin A derivatives used to raise antibodies. Cyclosporin A is linked to an immunogenic carrier via the amino acid side-chain at amino acid residue no. 1.\nCacalano, N. A.; Cleveland, W. L.; Erlanger, B. F. J. Immunol. Meth. 1989, 118, 257-263, describes the, presumably but not certainly, random photochemical grafting of cyclosporin A alkyl side-chains onto 4-benzoylbenzoic acid. These grafts are used to couple to immunogenic carriers for antibody preparation.\nU.S. Pat. No. 4,727,035 (1988), describes a method of immunoassay for cyclosporins. The method involves either radioiodine or fluorescence immunoassay.\nQuesniaux, V. F. J.; Tees, R.; Schreier, M. H.; Wenger, R. M.; Van Regenmortel, M. H. V. Molecular Immunology 1987, 24(11), 1159-1168, describes antibodies grown in response to cyclosporin-immunogen conjugates. The conjugates are attached to cyclosporin through a modified side chain on amino acids 2 or 8.\nU.S. Pat. Nos. 4,764,503 (1988), 4,639,434 (1987), and 4,384,996 (1983), describe cyclosporins modified at amino acid residue no. 8. These cyclosporins are hydroxylated on the amino acid side-chain and are useful as immunosuppressive drugs."}
-{"text": "A portion of the disclosure of this patent document contains material which is subject to copyright protection. The copyright owner has no objection to the facsimile reproduction by anyone of the patent disclosure, as it appears in the U.S. Patent and Trademark Office patent files or records, but otherwise reserves all copyright rights whatsoever.\nThe present invention relates to optical focusing systems and more particularly relates to a computer-based semi-automatic focusing and assembly apparatus and method for focusing and assembling optical devices, such as laser diode assemblies, used in bar-code scanning systems.\nBar-code scanning systems are widely used in retailing, postal and parcel delivery services, manufacturing, warehousing and distribution and other industries for fast data transaction applications. The benefits of such scanning systems are countless. For example, in manufacturing, the use of bar code scanners could result in accurate raw material inventory, production routing, work-in-progress tracking, labor efficiency, inventory control, quality assurance, shipping precision, and so on.\nOne of the key elements of bar-code scanners is the laser diode assembly that provides the light source for the scanners. The focusing characteristics of the assembly have significant effects on the correct reading of bar codes. If the laser beam emitted from the laser diode assembly is not focused accurately, it will result in a weak signal level and therefore possible loss of information when it scans across the bar-code. The end result is that one has to scan several times before correct information can be captured. For bar-code scanners to be effective over a range of distances, the beam width of the laser beam has to be within tolerance at several specified distances (normally three distances corresponding to the near, mid and far fields are specified). To ensure that the beam width requirement is met and therefore high scanner accuracy is achieved, every laser diode assembly needs to be focused in high precision during production.\nCurrently such focusing adjustments are performed manually in most of the bar-code scanner production processes which limits productivity as well as focusing performance of the laser diode assemblies. The manual focusing setup typically used in manufacturing laser diode (LD) assemblies consists of the following components: the laser diode assembly consisting of a light beam emitting laser diode, a focusing lens, a lens holder with an aperture, a spring, and a diode holder; a mechanical fixture to hold the LD assembly; a beam scan head including a cylindrical rotating drum with a slit for beam width measurement, a photo detector for beam power measurement, and a level for gain adjustment; a guide rail to move the scan head for beam width measurement at different distances (the height of the scan head can also be adjusted); a beam scan block which controls the scan head and displays the beam width (the outputs of the scan block include an analog beam profile signal, a trigger signal, and a clock reference signal); a current source that supplies the current to the laser diode; and a laser power meter which is used to measure the intensity of a laser beam during the laser diode power-up stage.\nThe first step of the focusing process is to assemble the laser diode, the spring, and the lens together in the diode holder and use a clamp to hold the parts in place. Next, this assembly is placed in the mechanical focusing fixture. The clamp is lifted and the assembly is now ready for focusing. For each type of assembly, the operator is usually given three specified locations at which the beam width should be focused (i.e., near, mid and far fields). The operator\"\"s next task is to connect the current source to the laser diode and slowly increase the current to a specified value. The operator then manually slides the scan head to the nearest specified location and adjusts the orientation of the LD so that the beam is centered in the aperture of the scan head. The operator then adjusts one or more focusing knobs of the focusing fixture so that the beam width displayed on the scan head is within a specified tolerance. The process is then repeated at mid and far field locations. If the operator is successful in achieving beam widths within the specified tolerances at all locations, glue is carefully applied to the lens holder to fix the lens with respect to the LD. After the curing time of the glue has elapsed, the operator removes the assembly and starts focusing the next assembly. If the operator is not able to focus the beam to acceptable beam widths at all locations, the assembly is rejected.\nFollowing are some of the major problems associated with the manual focusing setup and procedures:\nMoving Scan Head\nOne of the major concerns of the current setup is that the operator has to repeatedly move the scan head on the rail back and forth during focusing. For mass production, the time spent on moving the beam scan head significantly reduces productivity.\nTime Consuming Focusing Check After Curing\nAfter curing, the beam widths at specified distances have to be checked again to see if they are still within tolerances. To do this in the existing manual setup, the operator has to move the scan head to different locations again, which further slows down the process.\nHuman Factor\nIn the existing setup, the focusing criteria is quite subjective with the operator making most of the decisions based on intuition and experience. In addition, the operator is not required to optimize focusing.\nWasted Time During Curing\nCurrently, the operator focuses only one LD assembly at a time. During the curing period, the operator is basically idle. This period of time can be utilized more efficiently if the operator can start focusing the next assembly. With the current setup this is not possible.\nNo Check on Beam Centering\nAnother concern is that the operator cannot accurately point the laser beam to the center of the scan head aperture due to the lack of an accurate beam centering target. Because of this pointing error, the beam width displayed on the scan block may not be an accurate indication of the true beam width at the specified clip level.\nFocusing Inconsistency\nCurrently, the operator\"\"s task is simply to bring the beam widths within the specified tolerances. No targeting of specific beam width values within the specified ranges is required. The result of this practice is that the assembled LD assemblies lack focusing consistency from assembly to assembly, which in turn may result in performance inconsistency of bar-code scanner products.\nLack of Process Tracking\nDuring the focusing process, some of the focusing parameters may change systematically due to errors in the focusing setup or systematic defects in the laser diodes and focusing lenses. Since no focusing data are recorded in the current focusing process, no process tracking can be done to detect such problems. Also, if a laser diode assembly results in a defective product when assembled into a bar-code scanner, no trouble shooting can be done because no data regarding to the focusing condition of the laser diode assembly at the time of assembling is available.\nNo Control Over the Amount of Glue Dispensed\nAfter focusing, the amount of glue applied for curing and the way how glue is applied are crucial. The operator has to be very careful in making sure that no excessive amount of glue is applied. If this operation is not carried out with care, then the focal distance may change and the focusing may go out of tolerance. In the currently process, there is no control over the amount of glue dispensed, which could cause serious problems if the operator is not careful or not skillful.\nNo Access to Beam Profile Information\nAnother concern is that the engineers do not have easy access to the beam profile information at different distances, especially that for the assembled laser diode assemblies, to perform a comprehensive analysis of the process. The beam profiles together with beam widths can provide a more complete and in-depth knowledge about the nature of the laser beam.\nManual Power-up Procedure\nThe current procedure of powering up the LD is completely manual. The operator is given a target laser power level and a corresponding range of current values. The operator is supposed to manually adjust the current supply until the specified laser power level is reached. There is no specification nor control over the speed at which the operator should follow to power-up the LD. The major concern here is that rapid power-up and excessive current supply may result in permanent damage to the LD. Another concern is that there is no continuous monitoring of the laser power in the current focusing process. Therefore, if the laser power is changed, it cannot be detected and erroneous focusing may thus result.\nAccordingly, it is an object of the present invention to overcome such shortcomings and to provide an improved method and apparatus for focusing and assembling visible laser diode devices.\nThe present invention is an apparatus and method for semi-automatic focusing and assembling of a light beam modifier, such as an optical lens, relative to a light beam source, such as a laser diode, to optimize the light source\"\"s beam characteristics at one or more desired locations along the path of the beam. The present invention utilizes one or more light beam analyzers positioned at one or more desired locations along the path of the beam to measure a beam characteristic, such as beam width, and to generate value signals based on these measurements. The value signals are sent to a host computer for comparison with predetermined optimal beam parameter ranges. The computer generates a focusing signal, based on the comparison, which is used to axially adjust the light beam modifier relative to the light source.\nPreferably, the light beam modifier and light source are held in axial alignment in a focusing fixture which is adjustable for pointing the light beam in a desired direction and for focusing the beam. The focusing station includes a stepper motor electrically connected to the host computer for receiving the focusing signal and axially adjusting the light beam modifier based on this signal. The focusing station is mounted at a fixed position on a carrier rail arrangement. The one or more light beam analyzers are slidably mounted on the carrier rail arrangement for positioning the analyzers at desired locations.\nIn a preferred embodiment, the present invention includes two focusing stations, three light beam analyzers and two beam splitters mounted on a perpendicular carrier rail arrangement. The two focusing stations are mounted perpendicularly and equidistant from a first beam splitter. The first beam splitter splits the light beam emitted from either focusing stations into two components. The first component is directed to a first light beam analyzer positioned at a desired location and the second component is directed to a second beam splitter. The second beam splitter further splits from the second light beam component a third component. The second and third light beam components are respectively measured by a second and a third light beam analyzer positioned on the rails at desired locations.\nThe spatial arrangement of the focusing stations, two beam splitters and three light beam analyzers upon carrier rails in itself allows for simultaneous measuring of beam characteristics at three locations and for alternately focusing and assembling two laser diode assemblies without the aid of a host computer. Preferably, however, the three light beam analyzers generate value signals based on their respective beam measurements and send these signals to a host computer. The host computer receives these signals and compares them to predetermined parameter ranges associated with the selected locations along the beam path. The host computer generates a focusing signal based on an optimization algorithm for the three locations and sends the focusing signal to the stepper motors for axially adjusting the light beam modifiers with respect to the light beam sources.\nTo provide easy pointing of the laser beam emitted from either focusing station, the present invention also preferably includes a video camera electrically connected to the host computer and facing one of the scan heads. The video camera provides a live image of the laser spot on a computer screen along with a reference target. Pointing of the laser beam is done either by manually adjusting the pointing mechanisms of the focusing stations or automatically through the use of a data acquisition system of the computer. For automatic pointing, the computer preferably includes a video frame grabber which digitizes the laser beam image captured by the video camera and detects its location on a scan head. Again, an algorithm which can be provided by computer software or other signal generating means, e.g., analog and/or digital circuitry (printed or otherwise), generates an adjustment signal based on the difference between the detected beam location and the known center of the scan head. The adjustment signal is sent to additional adjustment motors which activate the pointing mechanisms of the focusing fixture to center the beam on the scan head.\nThe apparatus also preferably includes a semi-automatic glue dispensing system. A valve-controlled glue dispensing system with a custom-made glue dispensing gun is used. When the operator presses a switch on the glue dispensing gun, a well-controlled amount of glue is dispensed to fix the light beam modifier relative to the light source. Once this is done, a signal is sent to the computer via the data acquisition system which initiates a curing timer to count down. The timer alerts the operator when the proper curing time has elapsed.\nIn another preferred embodiment of the present invention, the apparatus includes automatic laser power-up circuitry. The computer first powers up the light source to a pre-focusing current level for pointing purposes. After pointing, the circuitry then powers the light source to the nominal power level. The power level is monitored continuously and the laser power is displayed on the computer screen.\nIn a method for optimizing characteristics of a light beam emitted from a light source at one or more locations along the path of the light beam according to the present invention, a light beam modifier is axially aligned with a light beam source and the beam emitted from the source is continuously measured at one or more selected locations along the path of the beam. Value signals corresponding to these measurements are generated and are compared with predetermined parameter ranges. A focusing signal based on the comparison is generated and is sent to the light beam modifier for axially adjusting the modifier with respect to the light source. The focusing signal terminates when the value signals fall within the predetermined parameter ranges and the light beam modifier is fixed relative to the light source. In a preferred method, two beam splitters are positioned along the path of the light beam to split the beam into three components which, in turn, are continuously measured at three selected locations.\nAs a result of the present invention, a method and apparatus is provided for semi-automatic focusing and assembling of laser diode assemblies is provided. The benefits of which are the reduction of the time needed to focus a laser diode assembly, the reduction of moving parts in the system as well as the overall improvement of quality. With the present invention, beam profiles along with other important focusing data are automatically saved to data files. These data files provide optical engineers with an extensive database of information for detailed analysis of the focusing process.\nFor a better understanding of the present invention, reference is made to the following detailed description to be taken in conjunction with the accompanying drawings and its scope will be defined in the appended claims."}
-{"text": "The present invention relates to a steel laminate gasket for an internal combustion engine to provide desired surface pressure around a hole to be sealed.\nAn engine for an automobile has been developed to seek high power and light weight. In order to satisfy these demands for the engine, an engine is made of aluminum alloy, to which high temperature and high pressure are applied.\nSince a steel laminate gasket is made of a plurality steel plates, the steel laminate gasket is strong against high temperature and pressure. Therefore, the steel laminate gasket has been used widely for a recent engine to which high temperature and pressure are applied.\nA cylinder head of a recent engine is generally made of aluminum alloy. In some cases, a cylinder block is also made of aluminum alloy. In case the cylinder head or cylinder block is made of aluminum alloy, a large amount of tightening pressure can not be applied to tighten the cylinder head and the cylinder block, because the cylinder head or cylinder block may deform by the tightening pressure.\nEspecially, when the gasket includes a wire ring around the cylinder hole, the cylinder head or cylinder block is liable to deform when tightened. Also, in case a large surface pressure is applied around a cylinder hole when the gasket is tightened, the cylinder head or cylinder block may deform around the cylinder hole.\nAs shown in FIGS. 1 and 2, a gasket 10 formed of an upper plate 11 with a bead 11a, a middle plate 12 and a lower plate 13 with a bead 13a is known. When the gasket 10 is tightened, the beads 11a, 13a deform to provide a strong surface pressure around a cylinder hole Hc. Therefore, unless the cylinder head and cylinder block are strong, the cylinder head and cylinder block may deform when tightened.\nIn U.S. Pat. No. 4,898,396, a gasket is formed of three plates, wherein upper and lower plates extend around a cylinder hole and have sealing mechanisms to seal therearound. In this patent, a surface pressure is formed adjacent to the cylinder hole.\nAccordingly, one object of the present invention is to provide a steel laminate gasket, in which the cylinder head and cylinder block can be tightened at desired surface pressure when the gasket is tightened.\nAnother object of the invention is to provide a steel laminate gasket as stated above, which can provide moderate surface pressure at an area slightly away from the cylinder hole to prevent deformation of the cylinder hole.\nA further object of the invention is to provide a steel laminate gasket as stated above, which can be easily and economically manufactured.\nFurther objects and advantages of the invention will be apparent from the following description of the invention."}
-{"text": "It is known to create three-dimensional models of physical surfaces, such as a forest, an urban area or the terrain of a utility corridor, from a location sampling of the physical surface obtained using a three-dimensional scanning range finder. For example, such measurements may be obtained by flying over the terrain of interest in an aircraft and scanning the ground surface using a scanning laser range finder such as \u201cLiDAR\u201d (Light Detection and Ranging).\nThe resulting data, once processed, will take the form of a group of spaced-apart, discrete, measured surface points (a \u201cpoint cloud\u201d) representing the physical surface in a three-dimensional coordinate reference frame (typically Cartesian). The spacing between the measured surface points results from the inherent limitations in resolution of the scanning range finder, and can be significant when the physical surface being scanned is large and the details of the physical surface are fine, such as when conducting an aerial survey of an urban area that is required to capture fine architectural details of individual buildings. Moreover, commercially available three-dimensional scanning range finders for terrain measurement only measure depth and not color. As a result, the visual representation generated from an airborne three-dimensional scan of a ground surface appears as a collection of uncolored or, at best, monochrome (gray level toned) dots, which can make it difficult to resolve details that are distinguished more by color than by shape. Although photographic images can be overlaid on models based on geometric solids that are derived from the surface samples, the derivation of such models is a lengthy process that requires considerable manual intervention. While there is an emerging technique known as \u201cflash LiDAR\u201d that obtains both range and color information, it is currently much more expensive than conventional LiDAR and does not provide sufficient surface detail or range accuracy for many applications."}
-{"text": "1. Field of the Invention\nThis invention relates to electrical connectors with a ground structure for cross talk control between signal carrying conductors and, in particular, where the connector is a right angle receptacle or a right angle header.\n2. Description of Related Art\nWith the advance of technology, a high density of electronic circuits and components can be located on a printed wiring board or printed circuit board (PCB). Along with this miniaturization of electronic circuits and components, electrical connectors are needed to electrically and mechanically interconnect one PCB, such as a back panel or mother board, to one or more other PCBs, such as daughter boards. Further, it is typically desirable for such connectors to have a high signal density capacity. That is, the connectors should permit a high number of signals to pass through the connector per unit volume of the connector. However, electrical signals carried on a conductor can interfere with a signal carried on an adjacent conductor.\nThis interfering electrical effect that an electrical signal carried on a given conductor exerts on a signal carried on an adjacent conductor is referred to as \"cross talk.\" Controlling this cross talk is especially important in high density connectors. Such control can be implemented in a variety of ways.\nOne method of controlling cross talk is to connect certain terminals in a high density connector to conductive areas of a printed circuit board that are in turn grounded or connected to a predetermined ground potential. This solution is external to the connector.\nU.S. Pat. No. 4,655,518 (to Lennart B. Johnson et al.), U.S. Pat. No. 4686,607 (to Lennart B. Johnson) and U.S. Pat. No. 4,869,677 (to Lennart B. Johnson et al.) disclose a daughter board/backplane assembly with contact elements dedicated for grounding purposes. Header contact elements have contacts that can be connected to ground or a predetermined potential on a backplane. The header contact elements have other spring contacts carried by an inside header wall for touching contacts carried by a right angle receptacle outer wall. Other contacts are integral with and perpendicular to the contacts carried by the right angle receptacle outer wall for connection to the daughter board.\nU.S. Pat. No. 4,601,527 issued to Timothy A. Lemke discloses an internal shielding structure for connectors, specifically in vertical and right angle headers. The shielding structure includes a ground strip affixed to a mating surface of a header housing. The shielding structure further includes an elongated conductive spring contact with contact beams that extend in holes of side walls of the housing, lock tabs that connect to the ground strip and ground bars for connection to a grounded chassis.\nU.S. Pat. No. 4,824,383 issued to Timothy A. Lemke discloses a shielding structure in connectors or plug-type terminators for either a multiple conductor cable or a multiple tracing substrate that electrically isolates individual or groups of contact elements in the terminator to prevent or minimize cross talk between adjacent conductors and to prevent or minimize degradation of signal transmission. The terminator includes a ground structure with generally U-shaped channels. Contact elements extend into the channels. The ground structure is connected to a predetermined potential, rather than dedicating some of the contact elements for this purpose.\nU.S. Pat. No. 4,898,546 issued to Richard A. Elco et al. discloses a ground shield device for right angle connectors. A different one of the shield devices straddles alternate columns of contact elements in the connector. Each shield device clips to a tail of one of the contact elements straddled by the shield device. The shield devices are connected to ground or a predetermined potential.\nIt is an object of this invention to provide a high density right angle electrical connector for electrically and mechanically interconnecting electronic circuits and/or components controlling cross talk within the connector.\nFurthermore, it is an object of this invention to provide a high density angled or right angle electrical receptacle for electrically and mechanically interconnecting a circuit assembly and a plurality of terminals arranged in rows and columns in a header or shroud to control impedance and/or cross talk thereby to reduce or minimize degradation of signal transmission within the receptacle.\nFurthermore, it is an object of this invention to provide a high density angled or right angle electrical header for electrically and mechanically interconnecting a circuit assembly and a plurality of terminals arranged in rows and columns in a receptacle to control impedance and/or cross talk thereby to reduce, prevent or minimize degradation of signal transmission within the header."}
-{"text": "It is quite known already to prepare acetylene polymers by using a catalytic system consisting of a compound of a transition metal and of an organometal compound of the elements of groups 1A, 2A, and 3A of the Periodic System.\nA preparation of this type is found described for the first time by G. NATTA et al. in 'Rend. Acc. Naz. Lincei, Scienze Fis. Mat. Nat.: 25, 3 (1958), according to which crystalline acetylene polymers are obtained in the form of a powder by using catalysts based on titanium and aluminum trialkyl alcoholates.\nMore particularly there is known the possibility of obtaining the polyacetylene directly in the polymerization in the form of a film (T. Iyo, H. Shirakawa & S. Ikeda, J. Polym. Sci. Polym. Chem., Ed. 13, 11 (1974).\nThis method is preferred in general even if with the same it is not possible to obtain films with good mechanical properties.\nIn the published Italian Patent Appl. No. 22722 A/82 there are described new polyacetylene films endowed with particular characteristics, such as the elongation at break, the breaking load and the molecular orientation under stretch. Such films are prepared by using as catalysts special titanium compounds, mainly dialcoxy-dicresoxy titanium in combination with aluminum trialkyls. In general, the values of apparent density of the polyacetylene prepared according to the known methods of the Prior Art prove to be comprised between 0.4 and 0.5 g/cc, with values of the surface area greater than 70 m.sup.2 /g.\nA considerable drawback of the polyacetylenes prepared so far is represented by their instability to oxidation, which considerably reduces the possibility of their being used, for instance, in the manufacture of electrodes for primary and secondary batteries."}
-{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device having a silicon-on-insulator (SOI) structure and a method for fabricating the same, and more particularly, to a method for fabricating a semiconductor device wherein a side wall oxide film or polysilicon layer is formed on the edge of an insulating film adapted to insulate the side wall of an active semiconductor substrate of the semiconductor device from a gate oxide film of the semiconductor device.\n2. Description of the Prior Art\nIn the fabrication of semiconductor devices, formation of a SOI structure is involved to achieve an isolation between adjacent elements, thereby obtaining a superior electrical characteristic. Such a SO structure is made by forming a silicon oxide film as an insulator on an under silicon substrate and forming another silicon substrate (to be used as an active substrate), for example, a single-crystalline silicon layer on the silicon oxide film.\nReferring to FIGS. 1A and 1B, a conventional SOI structure is illustrated.\nFIG. 1A is a view showing the layout of a metal oxide silicon field effect transistor (MOSFET) having the SOI structure. In FIG. 1A, respective positions of masks for an active region 1 and a gate electrode 2 on a silicon substrate (not shown) are shown.\nFIG. 1B is a cross-sectional view taken along the line I--I of FIG. 1A. A silicon oxide film 4 is deposited over a silicon substrate 3. A second silicon substrate 5 having a trapezoidal cross-sectional structure is formed on the silicon oxide film 4. On a desired portion of the resulting structure, a gate oxide film 6 and a gate electrode 7 are sequentially formed.\nAs shown in FIG. 1B, the second silicon substrate 5 has an inclined structure in that its thickness at the edge of the active region is smaller than its thickness d.sub.si at the middle portion of the active region. Due to such a structure, the depth of the depletion region is limited to the thickness of the second silicon substrate 5. As a result, the charge Q.sub.B of depleted bulk is limited by the thickness of the second silicon substrate 5, thereby decreasing the threshold voltage of the semiconductor device, as expressed the following equations: EQU V.sub.T =V.sub.FB +Q.sub.B /C.sub.OX (1) EQU C.sub.OX =.epsilon..sub.OX /t.sub.OX (2)\nwhere, V.sub.T represents the threshold, V.sub.FB the flat band voltage, Q.sub.B the bulk charge, C.sub.OX the capacitance of the oxide film, .epsilon..sub.OX the dielectric constant of the oxide film, and t.sub.OX the thickness of the gate oxide film.\nEquation (1) shows that it is impossible to increase the threshold voltage at the lower edge portion 8 of the second silicon substrate 5 even when the impurity doping rate at that portion increases. This implies a difference in threshold voltage between the middle portion and edge portion of the second silicon substrate 5. The drain current characteristic thus involves a point of inflexion when it varies in accordance with a variation in gate voltage, as shown in FIG. 2. As a result, it is difficult to control the threshold voltage of the semiconductor device having the SOI structure. Furthermore, the characteristic of the SOI element depends on the edge shape of the second silicon substrate 5. This results in a great variation in the characteristics of the semiconductor device."}
-{"text": "2-[4-(polyhalo-2-hydroxy-2-propyl)anilino]-2-oxazolines are known as disclosed in copending application Ser. No. 668,385, filed Mar. 19, 1976, now U.S. Pat. No. 4,081,547 as useful agents for the treatment of hypertension. The art also discloses that substituted anilino-2-oxazolines have hypotensive activity, e.g., U.S. Pat. Nos. 3,453,284; 3,499,083; and 3,499,084; and Belgian Pat. Nos. 704,392; 704,393; and 704,396. However, these compounds lack the 4-one or 2-thiazolin-4-one functional group present in the compounds of this invention."}
-{"text": "The present invention is directed to bipolar transistor logic gates, and more particularly to gates of this type that combine high speed operation with high packing density.\nLogic gates that are constructed utilizing Schottky transistor logic (STL) are noted for their high switching speeds, and find utility in a number of applications because of this characteristic. Basically, an STL gate comprises a Schottky transistor, i.e., a bipolar transistor with a Schottky diode clamp between its base and collector. The cathodes of a plurality of Schottky diodes are each connected to the collector of the Schottky transistor, and their anodes comprise the output terminals for the gate. In operation, the clamp between the base and collector of the Schottky transistor prevents this transistor from going into saturation. Consequently, the transistor is quick to respond to changes in the voltage levels at its base. Further information regarding the features and operating characteristics of an STL gate can be found in the article \"STL Technology\" by Ben J. Sloan, IEDM Technical Digest 1979, pp. 324-327.\nWhile STL gates are desirable from the standpoint of their high switching speed, one drawback that has heretofore been associated with them is the comparatively large area that is required for each gate, resulting in a low packing density on a chip. More particularly, in a conventional STL gate, each of the output Schottky diodes is formed on a separate island on the chip, which increases the area of the buried layer for the chip and adds to the isolation area on the chip The higher the fan-out for each gate, the greater the number of islands that must be provided per gate. For example, a conventional STL gate having a fan-out of four requires five islands, one for each of the four Schottky diodes and one for the transistor collector. The need to provide each of these separate areas on the chip limits the packing density that can be obtained."}
-{"text": "The present invention relates to benzazabicyclic carbamates of the formulae I and II below, and the pharmaceutically acceptable salts of such compounds. The compounds of formula I are cholinesterase inhibitors and are useful in enhancing memory in patients suffering from Dementia and Alzheimer's disease. The compounds of formula II are novel intermediates used in the synthesis of compounds of formula I. As set forth in detail below, the compounds of formula I and certain of the compounds of formula II are also useful as analgesic agents.\nAlzheimer's disease is associated with degeneration of cholinergic neurons in the basal forebrain that play a fundamental role in cognitive functions, including memory. Becker et al., Drug Development Research, 12, 163-195 (1988). As a result of such degeneration, patients suffering from the disease exhibit a marked reduction in acetylcholine synthesis, choline acetyltransferase activity, acetylcholinesterase activity and choline uptake.\nIt is known that acetylcholinesterase inhibitors are effective in enhancing cholinergic activity and useful in improving the memory of Alzheimer's patients. By inhibiting acetylcholinesterase enzyme, these compounds increase the level of acetylcholine, a neurotransmitter, in the brain and thus enhance memory. Becker et al., supra, report that behavioral changes following cholinesterase inhibition appear to coincide with predicted peak levels of acetylcholine in the brain. They also discuss the efficacy of the three known acetylcholinesterase inhibitors physostigmine, metrifonate, and tetrahydroaminoacridine.\nEuropean Patent 0253372 refers to 1,2,3,3a,8,8-ahexahydro -3a,8 (and 1,3a,8) -di (and tri) methylpyrrolo /2,3-b/indoles of the formula ##STR2## wherein R, R.sub.1, X, Z, and m are as defined in such patent, and states that such compounds inhibit acetylcholinesterase are useful as memory enhancing and analgesic agents.\nEuropean Patent 0154864 refers to physostigmine derivatives of the formula ##STR3## wherein R is (C.sub.2 -C.sub.20) alkyl, branched alkyl, cycloalkyl or aryl, and states that such compounds inhibit acetylcholinesterase and are useful in the treatment of Alzheimer's disease.\nYu et al., Febs Letters, 234, 1, 127-130, (1988), refer to physostigmine derivatives of the formula ##STR4## wherein R is as defined in such article, and discuss their relative potencies as inhibitors of acetylcholinesterase and butylcholinesterase as compared to the corresponding potency of physostigmine.\nAtack et al., J. Pharmacology and Experimental Therapeutics, 249, 1, 194-202 (1989), refer to certain carbamoyl and N(1)-substituted analogs of physostigmine and discuss their relative potency as cholinesterase inhibitors as compared to physostigmine.\nBrufani et al., Eur. J. Blochem, 157, 115-120 (1986), refer to physostigmine analogs of the formula ##STR5## wherein R.sup.1 is an alkyl group and R.sup.2 is hydrogen, and state that such compounds possess anticholinesterase activity.\nKnown cholinesterase inhibitors are useful over a relatively small range of concentrations and exhibit adverse side effects, becoming exremely toxic at concentrations substantially higher than the effective range. Also, the relationship between cholinesterase inhibition and changes in acetylcholine concentrations that follow such inhibition has been shown to be unpredictable, not solely the result of percent cholinesterase inhibition, and strongly affected by the properties of individual drugs. There is therefore a great need for novel cholinesterase inhibitors."}
-{"text": "A microcontroller unit (MCU) is small computer formed on an integrated circuit. MCUs can generate signals for controlling of a wide range of devices, such as electric motors, voltage regulators, office machines, appliances, implantable medical devices, etc.\nMCUs typically include a central processing unit (CPU), memory, and programmable peripherals components. The CPU executes a program (hereinafter referred to as an embedded program), which is typically stored in flash memory. MCUs may also include analog-to-digital converters (ADCs), digital-to-analog converters (DACs), comparators, timer/counter channels (hereinafter timer channels), etc. Timer channels are often used to autonomously control external devices such as induction motors. Because the timer channels operate autonomously, the CPU can perform other functions while timer channels control external devices."}
-{"text": "The invention relates to a process for the production of a vanadium/phosphorus mixed oxide catalyst precursor, its transformation into the active catalyst and a process for the production of maleic anhydride using this catalyst\nMaleic anhydride is a well known and versatile intermediate for manufacturing unsaturated polyester resins, pharmaceuticals or agrochemicals. It is usually produced by catalytic partial oxidation of aromatic (e. g., benzene) or non-aromatic (e. g., n-butane) hydrocarbons.\nThe main component of the active catalyst in the oxidation of non-aromatic hydrocarbons like n-butane to maleic anhydride is vanadyl pyrophosphate, (VO)2P2O7, which as a rule is obtained by thermal treatment of vanadyl acid orthophosphate hemihydrate of the formula (VO)HPO4xe2x80xa20.5H2O, acting as catalyst precursor. Both vanadyl pyrophosphate and vanadyl acid orthophosphate hemihydrate may, if desired, be accompanied by a promoter element selected from the groups IA, IB, IIA, IIB, IIIA, IIIIB, IVA, IVB, VA, VB, VIA, VIB and VIIIA of the periodic table of elements, or mixtures of such elements.\nMethods for preparing the precursor compound conventionally involve reducing a pentavalent vanadium compound under conditions which will provide vanadium in a tetravalent state (average oxidation number +IV).\nPrior art knows a great many different procedures, which however in general involve the use of vanadium pentoxide (V2O5) as a source of pentavalent vanadium and orthophosphoric acid (H3PO4) as the phosphorus source (see e.g. U.S. Pat. No. 5,137,860 or EP-A-0 804 963).\nAs a reducing agent in principle any inorganic or organic compound containing elements which are able to act as a redox couple possessing an oxidation potential suitable for this kind of reaction may be suitably applied.\nThe most common reducing agent is hydrogen chloride in aqueous solution.\nAlso favourably applied are organic media like primary or secondary aliphatic alcohols or aromatic alcohols such as benzyl alcohol as these compounds seem to at least in part dissolve the reactants and thereby facilitate the redox reaction.\nThe most preferred organic reducing agent is isobutyl alcohol as isobutyl alcohol combines optimal characteristics such as (i) a boiling point of 108xc2x0 C. at atmospheric pressure, (ii) dissolution of the vanadium alcoholates formed from V2O5, thus favouring a complete redox reaction in the liquid phase and (iii) achieving a redox potential for the couples isobutyl alcohol/isobutyraldehyde and isobutyl alcoholrisobutyric acid suitable to let the alcohol act as reducing agent. The tetravalent vanadium reacts with phosphoric acid (H3PO4) and leads to precipitation of the precursor vanadyl acid orthophosphate hernihydrate of the formula (VO)HPO4xe2x80xa20.5H2O. The precipitate is usually washed with isobutyl alcohol and subsequently dried.\nA major disadvantage of the conventional method as described above is that even after drying the precursor contains some percent of organic compounds from the organic reaction medium, compounds which are supposedly either (i) strongly adsorbed at the solid surface, and therefore not easily removable by the washing and drying treatment, or (ii) physically trapped in between the crystals of the precursor, or (iii) physically or chemically trapped a (xe2x80x9cintercalatedxe2x80x9d) in the crystalline structure of the precursor.\nIt has been found that this percentage of organic compound which remains trapped in the precursor is a fundamental parameter which can adversely affect the performance characteristics of the active catalyst obtained after the thermal treatment.\nThe object of the present invention therefore was to provide a method for controlling the carbon content in a vanadium/phosphorus mixed oxide catalyst precursor and accordingly to provide a superior catalyst precursor which, when activated, leads to superior results in the conversion of a non-aromatic hydrocarbon to maleic anhydride.\nIt was found that the objectives could be achieved with a new process for the preparation of a vanadium/phosphorus mixed oxide catalyst precursor according to claim 1.\nThe invention comprises reducing a source of vanadium in the presence of a phosphorus source in an organic medium which comprises\n(a) isobutyl alcohol or a mixture of isobutyl alcohol and benzyl alcohol and\n(b) a polyol\nin the weight ratio of 99:1 to 5:95.\nIn a mixture of isobutyl alcohol and benzyl alcohol the benzyl alcohol content is as a rule between 5 and 50 wt %.\nAs a source of vanadium a tetravalent or pentavalent vanadium compound may be applied. Representative examples, although not limiting, are vanadium tetrachloride (VCl4), vanadium oxytribromide (VOBr3), vanadium pentoxide (V2O3), vanadyl phosphate (VOPO4xe2x80xa2nH2O) and vanadium tetraoxide (V2O4). Vanadium pentoxide is the preferred vanadium source.\nAs mentioned above, the vanadium source may, if desired, be accompanied by promoter elements selected from the groups IA, IB, IIA, IIB, IIIA, IIIIB, IVA, IVB, VA, VB, VIA, VIB and VIIIA of the periodic table of elements, or mixtures thereof.\nPreferred promoter elements are selected from the group consisting of zirconium, bismuth, lithium, molybdenum, boron, zinc, titanium, iron and nickel.\nOrthophosphoric acid (H3PO4) is the preferred phosphorus source.\nIsobutyl alcohol is the preferred component (a).\nPolyols which can be used as the component (b) are expediently C2-6 aliphatic polyols, preferably C2-6-alkanediols such as 1,2-ethanediol, 1,2-propanediol, 1,3-propanediol, 1,2-butanediol, 1,3-butanediol, 1,4-butanediol, 2,3-butanediol, 1,2-pentanediol, 1,3-pentanediol, 1,4-pentanediol, 1,5-pentanediol, 2,3-pentanediol, 2,4-pentanediol, 1,2-hexanediol, 1,3-hexanediol, 1,4-hexanediol, 1,5-hexanediol, 1,6-hexanediol, 2,3-hexanediol, 2,4-hexanediol, 2,5-hexanediol and 3,4-hexanediol.\nMost preferred polyols are the C2-4-alkanediols 1,2-ethanediol, 1,3-propanediol and 1,4-butanediol.\nThe preferred mixture of alcohols contains 5 to 30 mol % of polyol with respect to isobutyl alcohol.\nAs a rule the vanadium source together with the phosphorus source is suspended in the organic medium and the mixture is kept under agitation at a temperature of expediently 90xc2x0 C. to 200xc2x0 C., preferably 100xc2x0 C. to 150xc2x0 C. over a period of 1 h to 24 h.\nThe ratio of vanadium source to phosphorus source is conveniently such that the P/V atomic ratio is in the range of 1:1 to 1.3:1, preferably 1.1:1 to 1.2:1.\nAs a rule after the reduction the precursor vanadyl acid orthophosphate hemihydrate of the formula (VO)HPO4.0.5H2O is formed which is filtered, washed and subsequently dried at a temperature of expediently 120xc2x0 C. to 200xc2x0 C.\nDue to the reduction treatment according to the invention the carbon content of the precursor can be controlled in the range of 0.7 wt. % to 15.0 wt. %, preferably in the range of 0.7 wt. % to 4 wt. %.\nIt has been found that best results are obtained with catalyst precursors which, after an additional thermal treatment at about 300xc2x0 C. for about 3 hours in air have a residual carbon content of 0.7 wt. % to 3 wt. %, most preferably 0.8 wt. % to 1.5 wt. %.\nOnce prepared the precursor can in view of its further activation treatment be formed into defined structures with defined properties. Such procedures may include wet grinding to a specific particle size, the addition of additives to improve attrition resistance, and the a formation of a convenient shape.\nA spherical shape for instance is most suitable for the application of the catalyst in a fluidized bed.\nThe further transformation of the so formed precursor into the active catalyst can be performed following a great number of activation processes known in the art, but in general include a heat treatment applying temperatures of up to 600xc2x0 C. More in detail, these processes may involve:\n(a) an initial heating of the precursor to a temperature not to exceed 250xc2x0 C.\n(b) a further heat treatment from about 200xc2x0 C. to at least 380xc2x0 C. to 600xc2x0 C. at the maximum\n(c) maintaining the temperature of stage (b) over a certain time and\n(d) cooling the activated catalyst, thereby maintaining an individually controlled atmosphere in all steps.\nIn a preferred embodiment, the activation of the precursor is accomplished using the procedure described in EP-A-0 804 963.\nIn a more preferred embodiment, the activation comprises the steps of\n(a) heating the catalyst precursor from room temperature to a precalcination temperature of about 300xc2x0 C. in air or oxygen-depleted air\n(b) keeping at said pracalcination temperature,\n(c) further heating the precalcined catalyst precursor in nitrogen up to a calcination temperature of about 550xc2x0 C. and\n(d) keeping at said calcination temperature.\nAfter the transformation into the active catalyst the catalyst is ready to be applied for the conversion of non-aromatic hydrocarbons to maleic anhydride.\nSuch processes are well known in the art, e. g. from U.S. Pat. No. 4,594,433, U.S. Pat. No. 5,137,860 or U.S. Pat. No. 4,668,652.\nIn general the non-aromatic hydrocarbon is converted with oxygen or an oxygen containing gas at a temperature from about 320xc2x0 C. to 500xc2x0 C. to maleic anhydride. The non-aromatic hydrocarbon is expediently selected from aliphatic C4-10 hydrocarbons, preferably nbutane. The conversion can take place in a fixed bed or a fluidized bed reactor but preferably in a fluidized bed reactor. The following examples are given by way of illustration only and arc not construed as to in any way limit the invention."}
-{"text": "The present invention relates to compounds with some structural similarities to bufadienolide compounds disclosed in the prior art. For a review of bufadienolides see Huimin Gao et al in Nat. Prod. Rep., 2011, 28, 953.\nThe bufadienolide compounds reported in the prior art are natural steroids, originally isolated from terrestrial natural sources such as plants of the families Crassulaceae, Hyacinthaceae, Iridaceae, Melianthaceae, Ranunculaceae, and Santalaceaethe, and animals of the genus Bufo (toads), Photinus (fireflies), and Rhabdophis (snakes) (Steyn et al. Nat. Prod. Rep. 1998, 15, 397-413; Krenn et al. Phytochemistry, 1998, 48(1), 1-29).\nAmong these bufadienolide compounds, Scilliroside and other Scilla compounds have been isolated from the red squill, Urginea maritima, and are disclosed to be highly toxic, specially Scilliroside which affects the cardiovascular and central nervous systems, causing convulsions and death (Verbiscar et al. J. Agric. Food Chem. 1986, 34, 973-979; Kopp et al. Phytochemistry, 1996, 42(2), 513-522). Majinda et al. also isolated bufadienolide compounds from Urginea sanguinea, which make the plant unsafe to be used as a medicinal plant (Planta Med. 1997, 63, 188-190).\nThe antiviral activity against a series of rhinoviruses and the anti-herpetic activity of several bufadienolide compounds were evaluated by Kamano et al. (Chem. Pharm. Bull. 1988, 36(1), 326-332) and Takechi et al. (Phytochemistry, 1996, 41(1), 125-127), respectively, and they found that most of the compounds displayed some inhibitory activity.\nAdditionally, the cytotoxic activity of several bufadienolide compounds has been evaluated by several authors. Specifically, Jing et al. reported that bufalin has a potent growth-inhibitory effect on human cancer cells of leukaemia (HL-60, ML1, U937, and K562 cell lines), epithelioid carcinoma (HeLa cell line), hepatoma (PLC/PRF/5 cell line), and epidermoid carcinoma (A431 cell line), but it is less potent on mouse leukaemia M1, melanoma B16, and lymphoid neoplasm P388 cell lines and rat hepatoma AH66 and chromaffin cell PC12 cell lines. They also found that bufalin induces typical apoptosis in human leukaemia HL-60 cell line but not in human leukocytes (Jpn. J. Cancer Res. 1994, 85(6), 645-651).\nKupchan et al. disclosed several bufadienolide compounds isolated from Bersama abyssinica which showed inhibitory activity against human carcinoma of the nasopharynx (KB) cell line (Bioorg. Chem. 1971, 1, 13-31; J. Org. Chem. 1971, 36(18), 2611-2616).\nKamano et al. evaluated the cytotoxic activity of 80 bufadienolide and cardenolide compounds, isolated from the Chinese drug Ch'an Su (obtained from the skin glands of toads such as Bufo gargarizans), against a primary liver carcinoma PLC/PRF/5 cell line and the colchicine-resistant cell line of PLC/PRF/5. Of them, 16 were shown to have potent cytotoxicities (IC50<10\u22123 \u03bcg/mL) against PLC/PRF/5 cell line (Bioorg. Med. Chem. 1998, 6, 1103-1115; J. Med. Chem. 2002, 45, 5440-5447). Additional bufadienolide compounds were isolated by Nogawa et al. from the same source, which were tested against human carcinoma of the nasopharyx (KB), human leukaemia (HL-60), murine leukaemia (MH60), pancreas adenocarcinoma (BXPC3), breast adenocarcinoma (MCF7), CNS glioblastoma (SF268), lung NSC (NCIH460), colon carcinoma (KM20L2), and prostate cancer (DU145) cell lines (J. Nat. Prod. 2001, 64, 1148-1152).\nYe et al. prepared novel bufadienolide compounds from bufalin by microbial hydroxylation. The compounds were tested against human hepatoma Bel-7402, human gastric cancer BGC-823, human cervical carcinoma HeLa, and human leukaemia HL-60 cell lines, showing some of them potent cytotoxicities comparable to that of bufalin (J. Steroid. Biochem. Mol. Biol. 2004, 91, 87-98; J. Nat. Prod. 2005, 68, 626-628).\nSince cancer is a leading cause of death in animals and humans, several efforts have been and are still being undertaken in order to obtain an anticancer therapy active and safe to be administered to patients suffering from a cancer. The problem to be solved by the present invention is to provide compounds that are useful in the treatment of cancer."}
-{"text": "Tremolo bridges having tremolo activation levers, which are often called whammy bars, vibrato bars, or tremolo bars, that are designed so that the musician can easily add vibrato while playing the guitar. Vibrato is the periodic variation in a tone and tremolo is periodic variation in the amplitude. Because of this confusion, tremolo and vibrato will be used interchangeably herein. This is accomplished by having the base plate of the bridge pivoting on a fixed point or points. The vibrato arm is connected to the base plate and when the musician moves the vibrato arm the strings of the guitar (or other instrument) increase and decrease in tension. As a result, the tone varies periodically resulting in a vibrato effect.\nOne problem with tremolo bridges is that when a musician pulls (bends) on a string, the tension of the other strings are affected. This occurs because the base plate is balanced on pivot point (line or anchor) and increasing the tension on one string results in base plate moving from its resting position. When the base plate moves the tension on the other strings change. In most cases it causes the tuning of the other strings to go slightly flat. This is an undesirable side effect of tremolo bridges and negatively effects the quality of musical sound.\nThere have been attempts to solve this problem by adding counterbalancing springs with different spring constants to the tremolo bridge. However, these have been ineffective since the tremolo bridge is already balanced between the force of the strings and the springs on the tremolo bridge. Adding additional springs does not solve the problem.\nThus there exists a need for a tremolo bridge that does not result in the other strings going flat when one of the strings is tensioned. Note that while the present invention is described with respect to guitars, it is applicable to any string instrument."}
-{"text": "1. Field of the Invention\nThe present invention relates to the conversion of thermal energy to electrical energy, and electrical energy to refrigeration, and more particularly to a solid state thermionic converter using semiconductor diode implementation.\n2. Relevant Technology\nThermionic energy conversion is a method of converting heat energy directly into electrical energy by thermionic emission. In this process, electrons are thermionically emitted from the surface of a metal by heating the metal and imparting sufficient energy to a portion of the electrons to overcome retarding forces at the surface of the metal in order to escape. Unlike most other conventional methods of generating electrical energy, thermionic conversion does not require either an intermediate form of energy or a working fluid, other than electrical charges, in order to change heat into electricity.\nIn its most elementary form, a conventional thermionic energy converter includes one electrode connected to a heat source, a second electrode connected to a heat sink and separated from the first electrode by an intervening space, leads connecting the electrodes to an electrical load, and an enclosure. The space in the enclosure is either highly evacuated or filled with a suitable rarefied vapor, such as cesium.\nThe essential process in a conventional thermionic converter is as follows. The heat source supplies heat at a sufficiently high temperature to one electrode, the emitter, from which electrons are thermionically evaporated into the evacuated or rarefied vapor-filled interelectrode space. The electrons move through this space toward the other electrode, the collector, which is kept at a low temperature near that of the heat sink. There the electrons condense and return to the hot electrode via external electrical leads and an electrical load connected between the emitter and the collector. The flow of electrons through the electrical load is sustained by the temperature difference between the electrodes. Thus, electrical work is delivered to the load.\nThermionic energy conversion is based on the concept that a low electron work function cathode in contact with a heat source will emit electrons. These electrons are absorbed by a cold, high work function anode, and they can flow back to the cathode through an external load where they perform useful work. Practical thermionic generators are limited by the work function of available metals or other materials that are used for the cathodes. Another important limitation is the space charge effect. The presence of charged electrons in the space between the cathode and anode will create an extra potential barrier which reduces the thermionic current. These limitations detrimentally affect the maximum current density, and thus present a major problem in developing large-scale thermionic converters.\nConventional thermionic converters are typically classified as vacuum converters or gas-filled converters. Vacuum converters have an evacuated medium between the electrodes. These converters have limited practical applications.\nEmbodiments in a first class of gas-filled converters are provided with a vaporized substance in the interelectrode space that generates positive ions. This vaporized substance is commonly a vaporized alkali metal such as cesium, potassium and rubidium. Because of the presence of these positive ions, liberated electrons can more easily travel from the emitter to the collector. The emitter temperature in this type of conventional device is in part determined by the vaporization temperature of the substance that generates the positive ions. Generally, the emitter temperature should be at least 3.5 times the temperature of the reservoir of the positive ion generating substance if efficient production of ions is to be achieved in these conventional devices.\nEmbodiments in a second class of gas-filled converters are provided with a third electrode to generate ions, and the gas in the interelectrode space in these conventional devices is an inert gas such as neon, argon, or xenon. Although these converters can operate at lower temperatures, such as about 1500 K, they are more complex.\nTypical conventional thermionic emitters are operated at temperatures ranging from about 1400 K to about 2200 K, and thermionic collectors are operated at temperatures ranging from about 500 K to about 1200 K. Under optimum conditions of operation, overall efficiencies of energy conversion range from 5% to 40%, electrical power densities are of the order of 1 to 100 watts/cm2, and current densities are of the order of 5 to 100 A/cm2. In general, the higher the emitter temperature, the higher the efficiency, the power and current densities with designs accounting for radiation losses. The voltage at which the power is delivered from one unit of a typical converter is 0.3 to 1.2 volts, i.e., about the same as that of an ordinary electrolytic cell. Thermionic systems with a high power rating frequently include many thermionic converter units connected electrically in series. Each thermionic converter unit is typically rated at 10 to 500 watts.\nThe high-temperature attributes of thermionic converters are advantageous for certain applications, but they are restrictive for others because the required emitter temperatures are generally beyond the practical capability of many conventional heat sources. In contrast, typical thermoelectric converters are operable at heat source temperatures ranging from about 500 K to about 1500 K. However, even under optimum conditions, overall efficiencies of thermoelectric energy converters only range from 3% to 10%, electrical power densities are normally less than a few watts/cm2, and current densities are of the order of 1 to 100 A/cm2.\nFrom a physics standpoint, thermoelectric devices are similar to thermionic devices. In both cases a temperature gradient is placed upon a metal or semiconductor, and both cases are based upon the concept that electron motion is electricity. However, the electron motion also carries energy. A forced current transports energy for both thermionic and thermoelectric devices. The main difference between thermoelectric and thermionic devices is in the transport mechanism; ballistic and diffusive transport for thermionics and ohmic transport for thermoelectrics. Ohmic flow is macroscopically diffusive, but not microscopically diffusive. The distinguishing feature is whether excess carriers are present. In thermoelectrics, the carriers normally present are responsible for current. In thermionics, the current is due to putting excess carriers in the gap. A thermionic device has a relatively high efficiency if the electrons ballistically go over and across the gap. For a thermionic device all of the kinetic energy is carried from one electrode to the other. The motion of electrons in a thermoelectric device is quasi-equilibrium and ohmic, and can be described in terms of a Seebeck coefficient, which is an equilibrium parameter.\nIn structures with narrow barriers, the electrons will not travel far enough to suffer collisions as they cross the barrier. Under these circumstances, the ballistic version of thermionic emission theory is a more accurate representation of the current transport.\nSolutions to the foregoing problems have been sought according to the present state of the art by using vacuum converters or gas-filled converters. Attempts to reduce space-charge effects with vacuum converters have involved the reduction of the interelectrode separation to the order of micrometers. Attempts to reduce the same effects with gas-filled converters have led to the introduction of positive ions into the cloud of electrons in front of the emitter. Nevertheless, these conventional devices still present shortcomings such as those related to limited maximum current densities and temperature regimes.\nConsequently, there remains a need to provide a more satisfactory solution to converting thermal energy to electrical energy at lower temperature regimes with high efficiencies and high power densities."}
-{"text": "The invention relates to the field of semiconductor packaging and testing, and in particular to mounting an integrated circuit that has been tested on a board, wherein the adjacent connecting pins of the integrated circuit are in an inline arrangement.\nIn order to keep the dimensions small for the package of a chip with multiple connecting pins (e.g., 88 pins), the connecting pins are arranged in line as close as possible. However, a minimum spacing must still be maintained since the space between associated holes on a board for insertion of the connecting pins may not be arbitrarily reduced. As the spacing of the holes and conducting lines on the board becomes closer, the process of fabricating a board becomes more complex and expensive. In addition, there is the risk that close spacing of the holes on the board will cause the connecting pins to short out during the soldering process.\nOne technique for increasing hole spacing on the board is to configure the connecting pins of the package by bending some of the connecting pins or, for example, by inwardly offsetting every second connecting pin. However, the advantage of greater hole spacing on the board is counteracted by the disadvantages of requiring a special test socket to test the chips and requiring an expensive packaging mechanism. Specifically, the packaging mechanism must be designed such that the connecting pins are not bent out of their precise alignment.\nTherefore, there is a need for a method of testing and mounting the package on a board in such a way that simple and inexpensively produced test sockets are sufficient for testing, and similarly simple and inexpensively produced packaging mechanisms are sufficient for packaging, while at the same time achieving an offset arrangement for the connecting pins."}
-{"text": "The present invention generally relates to perpendicular magnetic recording and reproducing systems, and more particularly to a perpendicular magnetic recording and reproducing system which realizes satisfactory perpendicular magnetic recording and reproducing characteristics.\nGenerally, when recording and reproducing a signal on and from a magnetic recording medium by use of a ring core head as a magnetic head, the ring core head magnetizes a magnetic layer of the magnetic recording medium in a longitudinal direction of the magnetic recording medium (that is, in an in-plane direction) at the time of the recording and picks up the recording at the time of the reproduction. However, according to such a longitudinal magnetic recording system, it is known that the demagnetization field becomes high as the recording density increases and the demagnetization field introduces undesirable effects to the high density magnetic recording. Hence, in order to eliminate the undesirable effects of the demagnetization, a perpendicular magnetic recording system has been proposed in which the ring core head magnetizes the magnetic layer of the magnetic recording medium in a direction perpendicular to the magnetic layer. According to the perpendicular magnetic recording system, the demagnetization field becomes low as the magnetic recording density increases, and theoretically, it is possible to realize a satisfactory high density magnetic recording in which there is no decrease in the remanent magnetization.\nAs a conventional perpendicular magnetic recording medium which is used in the perpendicular magnetic recording system, there is a perpendicular magnetic recording medium having a cobalt-chromium (Co--Cr) film formed on a base film by a sputtering process. As is well known, the Co--Cr film is extremely suited for use in the perpendicular magnetic recording medium because the Co--Cr film has a relatively high saturation magnetization (Ms) and favors magnetization in a direction perpendicular to the Co--Cr film (that is, the coercivity in the direction perpendicular to the Co--Cr film is large and the axis of easy magnetization is perpendicular to the Co--Cr film).\nHowever, when the ring core head performs the perpendicular magnetic recording and reproduction with respect to the perpendicular magnetic recording medium having the sputtered Co--Cr film, it is impossible to concentrate the magnetic flux at a predetermined magnetic recording position on the perpendicular magnetic recording medium, and there is a disadvantage in that it is impossible to obtain a strong magnetization which is in the direction perpendicular to the Co--Cr film and does not spread in the longitudinal direction of the perpendicular magnetic recording medium. In other words, when the ring core head is used to perform the recording on the Co--Cr film of the perpendicular magnetic recording medium, the magnetization direction easily deviates in the longitudinal direction of the perpendicular magnetic recording medium since the magnetic field generated by the ring core head includes considerable components in the in-plane direction. Accordingly, in order to maintain the magnetization direction in the perpendicular direction, the perpendicular magnetic recording medium must have a high perpendicular anisotropic magnetic field and have a saturation magnetization which is suppressed to a certain extent. However, the Co--Cr film does not have such characteristics, and there is a disadvantage in that it is impossible to perform a satisfactory perpendicular magnetic recording by the ring core head. In addition, the coercivity in the perpendicular direction must be large in order to obtain a high reproduced output from the perpendicular magnetic recording medium having the Co--Cr film. On the other hand, it is desirable to make the thickness of the perpendicular magnetic recording medium large in order to decrease the demagnetization field, but the perpendicular magnetic recording medium will not make contact with the ring core head in a satisfactory state when the thickness of the perpendicular magnetic recording medium is large because the perpendicular magnetic recording medium will lose its flexibility and become rigid. In this case, there are disadvantages in that the rigid perpendicular magnetic recording medium is easily damaged and undesirable effects are introduced to the ring core head, and it is impossible to perform a satisfactory perpendicular magnetic recording and reproduction.\nAccordingly, a perpendicular magnetic recording medium having a double film construction has been proposed. According to this perpendicular magnetic recording medium, a film having a high permeability, that is, a film having a low coercivity such as a nickel-ion (Ni--Fe) film, is formed between the Co--Cr film and the base film. The magnetic flux which is spread within the high permeability film is concentrated toward the magnetic pole of the ring core head at a predetermined magnetic recording position in order to obtain a strong magnetization which is in the perpendicular direction and does not spread in the longitudinal direction of the perpendicular magnetic recording medium. However, in the case of the perpendicular magnetic recording medium having the double film construction, the coercivity of the high permeability film is extremely small compared to the coercivity of the Co--Cr film, and there is a disadvantage in that Barkhausen noise is generated. For example, the coercivity of the Co--Cr film is over 700 Oe, and the coercivity of the high permeability film is under 10 Oe. Further, in order to produce the perpendicular magnetic recording medium having the double film construction, an amorphous (ion-nickel) Fe--Ni alloy or the like is formed on the base film by a sputtering process under a predetermined sputtering condition suited for forming the high permeability film, and Co--Cr is thereafter formed on the high permeability film by a sputtering process under a certain sputtering condition suited for forming the Co--Cr film. As a result, the sputtering condition under which the sputtering process is performed and the target must be changed for the formation of each film, and the sputtering processes cannot be performed continuously. Therefore, there are disadvantages in that the processes of manufacturing the perpendicular magnetic recording medium are complex and unsuited for mass production.\nOn the other hand, as the magnetic which is used in the perpendicular magnetic recording and reproducing system, there are in addition to the ring core head an auxiliary pole type (or driven) perpendicular magnetic head (hereinafter simply referred to as an auxiliary pole type head) having an auxiliary magnetic pole opposing a main magnetic pole and a single-sided main pole type (or driven) perpendicular magnetic head (hereinafter simply referred to as a single-sided main pole type head) not requiring an auxiliary magnetic pole and having a main magnetic pole opposing only one surface of the perpendicular magnetic recording medium. However, in the case of the perpendicular magnetic recording and reproducing system using the auxiliary pole type head, the perpendicular magnetic recording medium must be inserted between the main and auxiliary magnetic poles of the auxiliary pole type head, and the application of the system is limited. In other words, there is a disadvantage in that it is difficult to employ the auxiliary pole type head in a recording and reproducing apparatus such as a video tape recorder and a hard disc player. On the other hand, in the case of the single-sided main pole type head, the disadvantage of the auxiliary pole type head is eliminated, however, the construction of the single-sided main pole type head becomes complex. Furthermore, the single-sided main pole type head is disadvantageous in that the magnetic field distribution cannot be directed in the perpendicular direction as much as the auxiliary pole type head."}
-{"text": "Athletes engage in strenuous training to accomplish the goals of their sport. This strenuous training essentially amounts to trauma to the body, in that the human body interprets every strenuous work-out as a threat to its survival. It is known that muscle damage, caused by training, releases the catabolic hormone prostaglandin-E2. Training also causes the release of adrenocorticotropin (ACTH), which is a pituitary hormone. The presence of increased levels of ACTH increases the production of the catabolic hormone cortisol. Cortisol is also known as hydrocortisone, which is a glucocorticoid of the adrenal cortex that is a derivative of cortisone and is used in the treatment of rheumatoid arthritis. Thus, cortisol is a naturally occurring anti-inflammatory steroid. This catabolic hormone results in the release of amino acids from muscle tissue and prevents absorption of glucose. Cortisol, as a catabolic stress hormone, cannibalizes muscle tissue. High cortisol levels also result in the breakdown of connective tissue, lowered immunity and reduced muscle RNA synthesis. While cortisol may be a detriment to the athlete, scientists have conjectured that when the human body is stressed or traumatized, it triggers a \"fight or flight\" survival response. The biological design of cortisol is such that when a human is threatened, cortisol levels rise and mobilize the body for action by breaking down fat and muscle stores for emergency energy. Cortisol also reduces swelling in the event of injury. After the threat or trauma has subsided, cortisol levels return to normal. The cortisol-stress relationship is designed for intermittent physical threats and not the constant stimulation provided by today's aggressive athletes. Ongoing training results in cortisol levels that do not return to normal for extended periods of time and thereby result in the breakdown or loss of muscle tissue.\nAfter strenuous exercise, muscle tissue enters a stage of rapid nitrogen absorption in the form of amino acids and small peptides in order to rebuild the muscle fibers, grow and add new muscle fibers. During this period of repair and growth, it is important that the muscle cells have available to them sufficient levels of nitrogen in the form of amino acids. While the total level of amino nitrogen is important, the ratios of the various amino acids is even more important.\nAthletes that over-train sometimes enter into a catabolic condition. Muscle catabolism occurs when the athlete enters a negative nitrogen balance. People on diets usually have a negative nitrogen balance and therefore lose muscle when they lose weight. In contrast, a positive nitrogen balance means the animal has enough nitrogen left over to synthesize muscle proteins.\nVarious organizations have propounded a list of essential amino acids which are required on a daily basis for proper nutrition. These amino acid requirements vary throughout the growth cycle of all animals. Human muscle tissue is made up of specific amino acids and at specific ratios. If any of these amino acids are missing or deficient, the muscles will not grow, will grow slowly and may even begin to breakdown. However, if the animal is supplied with adequate amounts of protein that contain all of the muscle amino acids, this protein or source of amino nitrogen will be able to support rapid muscle recovery and growth.\nThe amino acids leucine, isoleucine and valine are the branched chain amino acids and are necessary for a positive nitrogen balance and muscle growth. The branched chain amino acids are lost at significant levels during strenuous exercise and therefore it is critical that they be available during the anabolic state.\nThe present invention is based in part, upon the discovery that the use of pyruvate in combination with an anabolic protein composition, produces a synergistic effect in increasing the lean body mass or muscle tissue of a mammal consuming same while at the same time decreasing the deposition of fat in the body.\nAs used herein and the claims, the term \"pyruvate\" means any salt or ester of pyruvic acid. Pyruvic acid has the formula: ##STR1##\nPyruvic acid is a colorless liquid with an odor resembling that of acetic acid and has a melting point of 13.degree. C. Pyruvic acid is an intermediate in the breakdown of sugars to alcohol by yeast. The mineral salts of pyruvic acid, such as magnesium pyruvate or calcium pyruvate or mixtures thereof are useful in the present invention. Sodium pyruvate is not especially preferred as it is known that sodium is associated with various negative medical conditions such as high blood pressure, water retention and heart disease. Further, certain athletes, such as bodybuilders, desire to present a defined visual image of their body which shows muscle definition and thus, the water retention properties of the sodium salt are not beneficial. Pyruvate precursors in the form of pyruvamides or pyruvyl-amino acids are also useful in the present invention. Pyruvyl-glycine is representative of the useful pyruvyl-amino acids.\nPyruvate has a number of useful applications in medicine. Pyruvate has been described for retarding fatty deposits in livers (U.S. Pat. No. 4,158,057); for treating diabetes (U.S. Pat. No. 4,874,790); for retarding weight gain (U.S. Pat. Nos. 4,812,879, 4,548,937, and 4,351,835); to increase body protein concentrations in a mammal (U.S. Pat. No. 4,415,576); for treating cardiac patients to increase the cardiac output without accompanying increase in cardiac oxygen demand (U.S. No. Pat. 5,294,641); for extending athletic endurance (U.S. Pat. No. 4,315,835); for retarding cholesterol increase (U.S. Pat. No. 5,134,162); for inhibiting growth and spread of malignancy and retarding DNA breaks (U.S. application Ser. No. 08/194,857 filed February 14, 1994); and for inhibiting generation of free radicals (U.S. application Ser. No. 08/286,946 filed Aug. 8, 1994). All of these references are incorporated herein by reference.\nU.S. Pat. 4,981,687 discloses compositions and methods for achieving improved physiological response to exercise. More specifically, this patent teaches a beverage comprising water, sugar, electrolytes, glycerol and pyruvate; and its use to ameliorate the effects of physical exertion. The teachings of U.S. Pat. No. 4,981,687 are incorporated herein by reference.\nU.S. Pat. No. 5,089,477 discloses the use of pyruvate in a liquid composition that is used to prevent weight loss in agricultural animals resulting from dehydration.\nU.S. Pat. No. 5,147,650 and U.S. Pat. No. 5,238,684 discloses and claims a fluid composition comprising water, electrolytes, sugar, glycerol, lactate and pyruvate. The teachings of U.S. Pat. Nos. 5,147,650 and 5,238,684 are incorporated herein by reference.\nU.S. Pat. No. 5,236,712 discloses and claims a beverage containing water, electrolytes, pyruvate and alanine in a concentration of from about 0.5% to about 10%. The teachings of U.S. Pat. No. 5,236,712 are incorporated herein by reference.\nPyruvate in various forms has been proposed for enteral administration and for parenteral administration. Typically, pyruvates are available in the form of salts, for example, calcium pyruvate and sodium pyruvate. U.S. Pat. Nos. 5,283,260 and 5,256,697 disclose uses for the pyruvyl-amino acids and methods for their production.\nPyruvate has been administered to mammals enterally or parenterally typically at superphysiological levels. The amount of pyruvate administered generally ranges from 1 to 20% of the mammal's caloric intake. For enteral dosage, the pyruvate may be disbursed or dissolved in a beverage product or may be included in cookies, candies or other foods. Pyruvate may also be introduced as an aqueous solution parenterally.\nThe amino acid composition of numerous proteins, biologically active polypeptides and foods, has been known for some time. One pre-eminent monograph on this subject is the Amino Acid Handbook by Richard J. Block, C. C. Thomas, publisher, 1956 (Library of Congress Catalog Card No. 56-9104). On page 343 at Table II, the approximate amino acid composition of mammalian muscle proteins on average, is set forth as follows:\n______________________________________ Amino Acid Grams/100 gm of Protein ______________________________________ Arginine 6.6 Histidine 2.8 Lysine 8.5 Tyrosine 3.1 Tryptophan 1.1 Phenylalanine 4.5 Cystine 1.4 Methionine 2.5 Serine 5.1 Threonine 4.6 Leucine 8.0 Isoleucine 4.7 Valine 5.5 Glutamic Acid 14.6 Aspartic Acid 8.0 Glycine 5.0 Alanine 6.5 Proline 5.0 Hydroxyproline 4.7 ______________________________________\nOn pages 272-273, the amino acid profile of human muscle tissue is set forth.\nProtein supplementation for serious athletes, such as body builders, is well accepted. Typically recommended dosages range between 2.0 and 3.5 gms of quality protein per kilogram of body weight per day. Numerous sources for the protein supplements are known, such as milk, egg, soy, beef and vegetable protein. Isolated fractions of these sources are also known such as ion-exchange whey protein, caseinates, whey protein concentrates, immunoglobulin and egg albumin. Protein supplements typically are provided as powders or tablets. It is also known to provide protein supplementation in the form of peptides (hydrolyzed protein) or even free amino acids. These approaches have two major limitations; cost and taste. While the use of pyruvate is known for various medical indications and the use of protein supplements are known to increase muscle mass, the prior art has failed to recognize or even consider the combination of pyruvate with an anabolic protein composition. The anabolic protein composition of the present invention is prepared through the combination of various proteins, peptides or amino acid sources to arrive at an amino acid profile that parallels the amino acid profile of human muscle tissue. Consumption of the anabolic protein composition will result in an increase in muscle mass, less catabolism of muscle tissue after strenuous exercise or as a result of disease, and quicker muscle cell restoration. Even more surprising are the synergistic results that can be achieved when pyruvate is combined with the anabolic protein composition of this invention. An additional aspect of the present invention is that the anabolic protein composition alone or in combination with pyruvate can lessen the catabolic effects of diseases such as cancer and AIDS. The prior art has not disclosed or suggested the compositions and methods of the present invention."}
-{"text": "During a manufacturing process of various structures, devices, systems, and the like, two separate and distinct components are secured together by one or more fasteners. The components may not have similar sizes or shapes. As an example, one of the components may be a flat planar piece, while the other component may be curved. In order to secure the two components together, a two piece washer assembly may be used. The washer assembly includes a first loose washer that receives a second loose washer. The two loose washers are configured to move through a limited range of motion in order to accommodate the angular differences between the two components.\nIn order to securely connect the two components together, an individual manipulates the washer assembly on one side of a first component, and a fastener on another side of the component. However, the washer assembly is susceptible to separating before receiving the fastener.\nAs can be appreciated, during the manufacture of a large system, such as an aircraft, the separate and distinct washers may be lost. For example, a flat drill plate may be secured to a portion of a curved fuselage panel. A first individual may be positioned on one side of the fuselage panel to secure the washer assembly, while a second individual may be positioned on an opposite side of the fuselage with a fastener. For large structures such as aircrafts, two individuals often work together to secure components (such as a drill plate and fuselage panel) together, through the use of the two piece washer assembly and a fastener. However, the unconnected loose washers may separate and fall into inaccessible areas of the fuselage. As can be appreciated, the manufacturing process may be labor and time intensive, and may be susceptible to loose parts being lost, or misplaced, such as within the system being manufactured."}
-{"text": "With increases in aging population, the development of pharmaceuticals effective for the treatment of senile dementia has been demanded strongly in recent years. Alzheimer's disease, which is a typical disease of senile dementia, is a neurodegenerative disease characterized by brain shrinkage, deposits of senile plaques, and neurofibrillary changes. This disease develops when insolubilized molecules attributed to a change of forms of amyloid \u03b2 peptides and following fibrillogenesis are deposited on neuronal cells such that neuronal cell death is induced by their toxicity.\nAmyloid \u03b2 peptides (A\u03b2) are degradation products generated by cleavage with \u03b2 secretase, etc., from neuronal amyloid precursor proteins, and these degradation products are found in 2 forms, A\u03b21-40 and A\u03b21-42. It has been reported that A\u03b21-42 is more likely to aggregate and more correlated to disease and neurotoxicity.\nIf a change of forms of amyloid \u03b2 peptides and fibrillogenesis thereof can be inhibited, the development of Alzheimer's disease can be prevented.\nThe pamphlet of WO 2005/105998 describes that a single-chain antibody having an activity of specifically binding to A\u03b21-42 such that the fibrillogenesis reaction thereof is inhibited is useful as a preventive or therapeutic agent for Alzheimer's disease.\nHowever, antibodies are high-molecular-weight proteins and are therefore expensive. They also have problems such as complicated production and purification steps or a lack of stability.\nDrugs having an effect of suppressing amyloid \u03b2 peptide production, for example, a rhodanine derivative (JP Patent Publication (Kokai) No. 06-192091A (1994)), a benzimidazole derivative (U.S. Pat. No. 5,552,426), a vinpocetine derivative (pamphlet of WO 96/25161), and an aromatic amide derivative (pamphlet of WO 2004/014843), have been known as low-molecular-weight preventive or therapeutic agents for Alzheimer's disease."}
-{"text": "The present disclosure relates generally to information handling systems, and more particularly to programming fuses in chipsets used in information handling systems.\nAs the value and use of information continues to increase, individuals and businesses seek additional ways to process and store information. One option available to users is information handling systems. An information handling system generally processes, compiles, stores, and/or communicates information or data for business, personal, or other purposes thereby allowing users to take advantage of the value of the information. Because technology and information handling needs and requirements vary between different users or applications, information handling systems may also vary regarding what information is handled, how the information is handled, how much information is processed, stored, or communicated, and how quickly and efficiently the information may be processed, stored, or communicated. The variations in information handling systems allow for information handling systems to be general or configured for a specific user or specific use such as financial transaction processing, airline reservations, enterprise data storage, or global communications. In addition, information handling systems may include a variety of hardware and software components that may be configured to process, store, and communicate information and may include one or more computer systems, data storage systems, and networking systems.\nMany information handling systems include chipsets such as, for example, the Platform Controller Hub (PCH) available from INTEL\u00ae Corporation of Santa Clara, Calif., United States. Such chipsets may include One-Time-Programmable (OTP) Non-Volatile Memory (NVM) and/or other memory subsystems that utilize fuses such as In Field Programmable (IFP) fuses that may be burned during manufacture to provide security information in the chipset (e.g., information associated with a public key such as a hash of a master public key) and/or enable particular features in the system. For example, security information provided in the chipset in such a manner and associated with a public key may be utilized to verify software or firmware (e.g., a Basic Input/Output System (BIOS)) that has been signed with an associated private key. Similarly, features enabled in the information handling system in such a manner may include security features such as, for example, BOOT GUARD available in systems provided by DELL\u00ae Inc. of Round Rock, Tex., United States, and PLATFORM TRUST TECHNOLOGY available from INTEL\u00ae corporation of Santa Clara, Calif., United States. Providing chipsets with information and enabling system features in such manner provides a root of trust for that information and those features.\nConventional systems provide two methods for programming IFP fuses during manufacture: Automatic End Of Manufacturing (Auto EOM) and Host Embedded Controller Interface (HECI) EOM. Auto EOM provides for the burning of fuses in the chipset automatically upon the initial boot of the system and before the BIOS begins executing. As such, Auto EOM cannot be controlled by the BIOS or any software tool or script. HECI EOM provides for the burning of the fuses in response to a command provided post-boot (e.g., using an SPSManuf Universally Extensible Firmware Interface (UEFI) shell tool or executable version that can be run in an operating system.) Neither of Auto EOM or HECI EOM differentiates between the chipset (e.g., the PCH) or the Central Processing Unit (CPU) installed in the system, which can raise some issues. For example, each of PCH and CPU may be provided in different versions (e.g., a production version (also referred to as a Qualification Sample (QS)) or a pre-production version (also referred to as an Engineering Sample (ES)). The IFP fuse burn operations may differ based on the combination of the PCH and the CPU utilized in a system, and conventional BIOS storage only stores a single set of IFP fuse burn instructions. As such, a manufacturer must determine and provide the appropriate IFP fuse burn instructions in the BIOS storage for use in the IFP fuse burning process, and risks rendering the system unusable (e.g., \u201cbricked\u201d) if the wrong IFP fuse burn profile is utilized.\nAccordingly, it would be desirable to provide an improved chipset fuse programming system."}
-{"text": "Ovarian cancer is a significant health problem for women in the United States and throughout the world. Although advances have been made in detection and therapy of this cancer, no vaccine or other universally successful method for prevention or treatment is currently available. Management of the disease currently relies on a combination of early diagnosis and aggressive treatment, which may include one or more of a variety of treatments such as surgery, radiotherapy, chemotherapy and hormone therapy. The course of treatment for a particular cancer is often selected based on a variety of prognostic parameters, including an analysis of specific tumor markers. However, the use of established markers often leads to a result that is difficult to interpret, and high mortality continues to be observed in many cancer patients.\nImmunotherapies have the potential to substantially improve cancer treatment and survival. Such therapies may involve the generation or enhancement of an immune response to an ovarian carcinoma antigen. However, to date, relatively few ovarian carcinoma antigens are known and the generation of an immune response against such antigens has not been shown to be therapeutically beneficial.\nAccordingly, there is a need in the art for improved methods for identifying ovarian tumor antigens and for using such antigens in the therapy of ovarian cancer. The present invention fulfills these needs and further provides other related advantages."}
-{"text": "The invention relates to an intermediate plate for mounting between housing parts of a fluid-operated control unit, in particular a gear unit, including a central unit that is constructed in the shape of a plate and, on either side of the central unit, a respective sealing system for sealing between the central unit and the housing part opposite the latter, this sealing system including sealing elements which are arranged on the respective side of the central unit and provide sealing around passages.\nIntermediate plates of this kind are known from the prior art.\nIn these, the problem arises of achieving the best possible seal between the central unit and the respective housing part that also maintains leak-tightness in the event of any kind of deformation of the housing parts and the central unit, even in the event of the high fluid pressures that are conventionally present in the case of a gear unit."}
-{"text": "This invention relates to a cutting tool guiding system for metal cutting machines.\nIn the manufacture of steel plates and sheets, it is desirable to start with metal slabs which are an optimum width for the particular end product desired. Many such slabs are formed on continuous casting machines and are often wider than necessary. It is usually the practice to slit such slabs longitudinally into two or more strips.\nThe original slabs are often separated from the continuous cast strand by a flame cut-off device which utilizes a cutting torch which is moved transversely across the strand. The separated slab has often been found to have corners which are deformed in that they contain slag, or are rounded. These deformities have caused problems with prior guidance systems for slitting the slab when such systems followed the slab edge. Furthermore, at least some prior edge following devices have also had to contend with warped or misplaced slabs.\nThe present invention overcomes the aforementioned problems and provides a longitudinal cutting device which is edge-guided by a system which is not affected by corner deformities or warpage or misplacement of the slab.\nIn accordance with the various aspects of the invention, a cantilevered transverse arm is mounted for movement along both a longitudinal and transverse axis. The outer end of the arm is adapted to be disposed over a metal slab to be cut and carries a tool assembly which is adjustable thereon. A guiding assembly is fixedly mounted to the arm and includes a pair of longitudinally spaced guide rollers positioned so that the cutting tool is disposed in a transverse plane extending between the rollers. Each roller is individually shiftable into and out of engagement with the slab edge.\nInitially, the leading roller is advanced toward the slab edge, while the trailing roller is retracted. The rollers are then moved toward the slab edge until the leading roller engages the said edge ahead of the slab corner. As the entire device moves longitudinally along the slab, the leading roller is retracted while the trailing roller moves into engagement with the slab edge. Deformed slab corners are thereby bypassed by the tool guidance system."}
-{"text": "The present invention relates generally to the field of touch screen devices, and more particularly to selecting multiple areas of content on a touch screen device.\nA touch screen is an input device normally layered on the top of an electronic visual display of an information processing system. A user can give input or control the information processing system through simple or multi-touch gestures by touching the screen with a special stylus/pen and/or one or more fingers. Some touch screens use ordinary or specially coated gloves to work while others use a special stylus/pen only. The user can use the touch screen to react to what is displayed and to control how it is displayed (for example by zooming the text size). The touch screen enables the user to interact directly with what is displayed, rather than using a mouse, touchpad, or any other intermediate device (other than a stylus, which is optional for most modern touch screens).\nTouch screens are common in devices such as game consoles, personal computers, tablet computers, and smartphones. They may also be attached to computers or, as terminals, to networks. They also play a prominent role in the design of digital appliances such as personal digital assistants (PDAs), global positioning system (GPS) navigation devices, mobile phones, video games and some books (electronic books or E-books). The popularity of smartphones, tablets, and many types of information appliances is driving the demand and acceptance of common touch screens for portable and functional electronics. Touch screens are found in the medical field and in heavy industry, as well as for automated teller machines (ATMs), and kiosks such as museum displays or room automation, where keyboard and mouse systems do not allow a suitably intuitive, rapid, or accurate interaction by the user with the display's content."}
-{"text": "This section is intended to provide background information to facilitate a better understanding of various technologies described herein. As the section's title implies, this is a discussion of related art. That such art is related in no way implies that it is prior art. The related art may or may not be prior art. It should therefore be understood that the statements in this section are to be read in this light, and not as admissions of prior art.\nA radar may provide marine vessels with the capability to safely navigate a course through unknown bodies of water. As such, a radar may prove to be a useful tool in navigation by identifying neighboring vessels and obstacles in close proximity to a marine vessel."}
-{"text": "The present invention relates to entry rug cleaning assemblies which remove loose soil from the soles of the shoes of persons passing thereover and in which the surface of the rug is cleaned periodically.\nIn applicant's prior U.S. Pat. No. 3,526,015 granted Sept. 1, 1970, there is disclosed an entry rug cleaning mechanism which has enjoyed commercial acceptance and which has proven quite useful in cleaning and disinfecting the soles of shoes of persons passing thereover into a building. In this installation, the rug is mounted so that it may be rotated about a closed loop, and within an adjacent housing disposed below the traffic surface are vacuum cleaning and washing installations for cleaning the surface of the rug passing thereover.\nSelf cleaning entry rug installations of this type are extremely desirable from the standpoint that fresh rug surface may be provided to effectively remove soil from the soles of the persons passing thereover and they are able to provide a relatively clean rug surface from an aesthetic standpoint as compared with fixed rug installations for this purpose.\nAlthough the installations of applicant's prior patent have provided substantial benefits, location of the cleaning components outside the belt loop increases the space required for the installation and does not provide optimum cleaning action since the area available for cleaning action is always a vertical portion of the belt loop. Moreover, frequent servicing of the cleaning components has been required.\nIn addition, replacement of the rug was frequently difficult and time consuming, and the tendency of the rug to stretch required constant adjustment of the tension.\nAccordingly, it is an object of the present invention to provide a novel self-cleaning entry rug assembly providing enhanced cleaning action and adapted for facile servicing and relatively trouble free operation.\nIt is also an object to provide such an entry rug assembly in which the components may be readily fabricated and assembled from relatively durable elements to provide a rugged installation.\nAnother object is to provide such an entry rug assembly which is adapted to fully automatic control, semi-automatic control, or manual operation.\nStill another object is to provide such an assembly which may be linked to presence sensors or approach sensors to limit operation during the time when persons are passing thereover."}
-{"text": "The invention relates to a composite driving belt provided with a carrier and a plurality of transverse elements assembled slidably thereon."}
-{"text": "1. Field of the Invention\nThe present invention relates to surgical devices and methods to guide instruments that prepare the surface of bones and other tissues for implants that replace a damaged, diseased, or otherwise painful spinal facet joint.\n2. Description of Related Art\nTraumatic, inflammatory, metabolic, and degenerative disorders of the spine can produce debilitating pain that can have severe socioeconomic and psychological effects. One of the most common surgical interventions today is arthrodesis, or spine fusion, of one or more motion segments, with approximately 300,000 procedures performed annually in the United States. Clinical success varies considerably, depending upon technique and indications, and consideration must be given to the concomitant risks and complications. For example, Tsantrizos and Nibu have shown that spine fusion decreases function by limiting the range of motion for patients in flexion, extension, rotation, and lateral bending. Furthermore, Khoo and Nagata have shown that spine fusion creates increased stresses and, therefore, accelerated degeneration of adjacent non-fused motion segments. Additionally, pseudoarthrosis, as a result of an incomplete or ineffective fusion, may reduce or even eliminate the desired pain relief for the patient. Finally, the fusion device, whether artificial or biological, may migrate out of the fusion site.\nRecently, several attempts have been made to recreate the natural biomechanics of the spine by use of an artificial disc. Artificial discs provide for articulation between vertebral bodies to recreate the full range of motion allowed by the elastic properties of the natural intervertebral disc that directly connects two opposed vertebral bodies.\nHowever, the artificial discs proposed to date do not fully address the mechanics of motion of the spinal column. In addition to the intervertebral disc, posterior elements called the facet joints help to support axial, torsional and shear loads that act on the spinal column. Furthermore, the facet joints are diarthroidal joints that provide both sliding articulation and load transmission features. The effects of their absence as a result of facetectomy was observed by Goh to produce significant decreases in the stiffness of the spinal column in all planes of motion: flexion and extension, lateral bending, and rotation. Furthermore, contraindications for artificial discs include arthritic facet joints, absent facet joints, severe facet joint tropism or otherwise deformed facet joints, as noted by Lemaire.\nU.S. Pat. No. Re. 36,758 to Fitz discloses an artificial facet joint where the inferior facet, the mating superior facet, or both, are resurfaced.\nU.S. Pat. No. 6,132,464 to Martin discloses a spinal facet joint prosthesis that is supported on the posterior arch of the vertebra. Extending from this support structure are inferior and/or superior blades that replace the cartilage at the facet joint. The Martin prosthesis generally preserves existing bony structures and therefore does not address pathologies that affect the bone of the facets in addition to affecting the associated cartilage. Furthermore, the Martin invention requires a mating condition between the prosthesis and the posterior arch (also known as the lamina) that is a thin base of curved bone that carries all four facets and the spinous process. Since the posterior arch is a very complex and highly variable anatomic surface, it would be very difficult to design a prosthesis that provides reproducible positioning to correctly locate the cartilage-replacing blades for the facet joints.\nAnother approach to surgical intervention for spinal facets is provided in WO9848717A1 to Villaret. While Villaret teaches the replacement of spine facets, the replacement is interlocked in a manner to immobilize the joint.\nIt would therefore be an improvement in the art to provide a vertebral facet replacement device and method that provides a relatively high degree of mobility in the joint, while effectively removing the source of arthritic, traumatic, or other disease mediated pain with a minimum of patient discomfort."}
-{"text": "1. Field of the Invention\nThe present invention concerns a solenoid valve with programmable electronic control device, particularly for watering systems.\n2. Description of the Related Art\nThere are also known solenoid valves of the aforesaid type in which the control solenoid is made up of two parts, as described in U.S. Pat. No. 6,352,238, where a first part of the solenoid includes a piston for the operation of a hydraulic valve and a magnet for attraction of the same piston and is screwable on the external body of the solenoid valve and a second part includes an electric solenoid for the excitation of the magnet that is provided with xe2x80x9cfaston(copyright)xe2x80x9d type electric connectors and which is fastened to the first part in a removable way by means of a removable and re-insertable locking element.\nFor some uses it can be desirable to equip the single solenoid valve with its own electronic control device of programmable type, that makes the same solenoid valve independent from the general program of the watering system or that even prevents the need of having a control unit for the system.\nHowever a solenoid valve provided with an independent control device inevitably becomes object of thefts and vandalism acts, particularly during the periods of absence of the user.\nIn addition the electronic control device inset in the solenoid valve can deteriorate for infiltration of humidity or rain or other atmospheric events, in particular during the autumn and winter periods, when besides the watering is no longer required.\nFinally the programming operation of the electronic control device inset in the solenoid valve requires the intervention of the user or the programmer directly on the site where it is used, often inconvenient and difficultly accessible.\nAll that considered, object of the present invention is to realise a solenoid valve with autonomous programmable electronic control device that is free of the aforementioned disadvantages.\nAccording to the present invention, such object is reached by means of a solenoid valve comprising a hydraulic valve, a valve body in which said hydraulic valve is housed and a control solenoid for said valve, characterised in that it comprises in addition a programmable electronic control device that is housed together with said solenoid in a removable way inside a watertight openable container that is sealingly fastened to said valve body.\nOwing to the present invention a solenoid valve with its own electronic control device is thus realised that, when in use, preserves the same device as regards infiltration of humidity or rain or others atmospheric events and at the same time it allows, when desired, to extract the control device so as to have it protected from thefts and vandalism acts and to allow its convenient programming at home or in a laboratory.\nPreferably the solenoid is of the type that is made up of a first part including a piston for the operation of the hydraulic valve and an attraction magnet for said piston and of a second part fastened to the first one in a removable way and including an electric solenoid for the excitation of the magnet.\nThe electric solenoid too, that it is the most expensive part and subject to failure of the solenoid, can therefore take advantage of the favourable effects of the watertight containment and removability as the ones previously described for the electronic control device."}
-{"text": "It is known in the art to use dental appliances for the treatment of snoring and sleep apnea. For example and without being limitative, the dental appliances can be intraoral orthoses which prevent the mandible from moving posteriorly during sleep.\nOne type of known intraoral orthosis includes two dental trays (or bite forms) joined together so that the user's mandible projects forwardly from its normal position and thereby helps reducing snoring and sleep apnea. In such an orthosis, the dental trays are connected to one another through two lateral links which are configured to advance the user's mandible during a vertical movement between the maxilla and the mandible. In some embodiments, for connection to the dental trays corresponding to the maxilla or the mandible, the lateral links have a protrusion inserted in a through hole of the dental trays.\nOne of the common drawbacks of known intraoral orthoses is that the lateral links often disconnect from the dental tray to which they are engaged, when the intraoral orthosis is being manipulated. In such cases, the user must spend time to reconnect the disconnected lateral link to its respective dental tray.\nIn view of the above, there is a need for an improved mandibular protrusion device which, by virtue of its design and components, would be able to overcome or at least minimize some of the above-discussed prior art concerns."}
-{"text": "1. Field of the Invention\nThis invention generally relates to a cycle computer display apparatus for a bicycle. More specifically, the present invention relates to a cycle computer display apparatus having a touch screen.\n2. Background Information\nBicycling is becoming an increasingly more popular form of recreation as well as a means of transportation. Moreover, bicycling has become a very popular competitive sport for both amateurs and professionals. Whether the bicycle is used for recreation, transportation or competition, the bicycle industry is constantly improving the various components of the bicycle.\nRecently, bicycles have been equipped with electrical components to make riding easier and more enjoyable for the rider. Some bicycles are equipped with automatic shifting units that are automatically adjusted according to the riding conditions by a cycle computer or control unit. In particular, the front and rear derailleurs have recently been automated. Moreover, various electronic devices have been used to determine one or more operating parameters for controlling the derailleurs and providing information to the rider. Thus, the cycle computer or control unit of the bicycle not only provides information to the rider, but also controls various components of the bicycle.\nIn view of the above, it will be apparent to those skilled in the art from this disclosure that there exists a need for an improved cycle computer display apparatus. This invention addresses this need in the art as well as other needs, which will become apparent to those skilled in the art from this disclosure."}
-{"text": "The process known as electrophoresis involves the migration of charged molecules through a suitable retarding medium under the influence of an electric field. Generally, the compounds of higher molecular weight migrate at a slower rate through the medium than do the compounds of lower molecular weight. Devices have been provided previously for carrying out electrophoresis. An example of such a device is U.S. Pat. No. 4,415,418 in which a tray is provided with a raised platform at the center. Removable partitions are placed in the tray at opposite ends of the platform, and a conventional electrophoresis gel is poured over the platform to form a thin layer. When the gel has cooled, the partitions are removed. A comb is provided to form wells across the surface of the gel. Substances that are to be subjected to electrophoresis are delivered into each of the wells, and the tray is at least partially filled with an electrolyte buffer. Electrodes are positioned at each end of the tray and a sufficient voltage difference is applied to the electrodes to cause migration of the molecules of the substance in the wells across the length of the gel, separated according to their molecular weight. After electrophoresis, the gel is removed from the original casting tray, and placed in a dish containing depurination solution. Approximately thirty minutes later this solution is poured out by tipping the dish toward one edge while the gel is held with the fingers. It is important to use great care during this procedure to prevent the gel from breaking because there is no gel support structure and subsequent processing is possible only with an integral gel. A denaturation solution is then added to the dish and incubation is continued for approximately thirty minutes. Again, the solution is carefully poured off. Then neutralization buffer is added and incubation is continued for thirty additional minutes.\nIn accordance with conventional techniques, transfer of the nucleic acids is accomplished by placing a piece of filter paper, which is as wide as and longer than the gel, on a platform which is suspended above a solution of 10.times. saturated saline citrate buffer (SSC). The ends of the filter paper are long enough to hang off the ends of the platform and dip into the 10.times. SSC. Thus, the filter paper acts as a wick to absorb the SSC solution. The gel is removed from the dish and placed on top of the filter paper saturated with 10.times. SSC. Next, a piece of membrane filter paper which is the same size as the gel is saturated with 10.times. SSC and placed on top of the gel. The nucleic acids are eventually bound to the membrane filter paper. Another piece of saturated filter paper, the same size as the gel, is placed on top of the membrane. The entire layered unit is then smoothed to remove any air bubbles that may exist between the gel and the filter paper. Finally, a stack of paper towels, the same size as the gel, is positioned on top of the layered unit.\nOver a period of about 12 to 16 hours, the 10.times. SSC solution is drawn up through the gel by capillary action and the nucleic acids are transferred out of the gel into the membrane above. The paper towels absorb the excess buffer and provide the force for capillary action. At the end of the transfer period, the entire layered unit is disassembled and the membrane is removed for hybridization. This technique is described in an article by E. Southern, \"Detection of Specific Sequences Among DNA Fragments Separated by Gel Electrophoresis,\" J. Mol. Biol., 98:503 (1975).\nAlthough the trays such as the one described in U.S. Pat. No. 4,415,418 are convenient for carrying out electrophoresis, they are not suitable for situations where a large number of samples must be tested in a relatively short period of time.\nTherefore, the prior art uses a tedious multi-step, multi-apparatus process for preparing nucleic acid fragments for subsequent hybridization. Four steps are generally undertaken to achieve preparation of the sample for hybridization. Electrophoresis was previously described. Depurination removes purine bases from nucleic acids. Denaturation involves separating the strands of nucleic acids and breaks down the depurinated nucleic acids into suitable size to allow eventual transfer of the fragments out of the gel. Transfer involves allowing the fragments to go out of the gel onto the porous membrane."}
-{"text": "Multi-lumen catheters are used for the purpose of creating two or more separate fluid pathways, such as in hemodialysis applications. A primary goal of hemodialysis access is to provide a reliable and effective means of dialysis, which means that a sufficient volume of blood over a period of time must be removed from and returned to the patient. Because the contaminated and cleansed blood must be kept separate for an effective dialysis procedure, a dual lumen catheter is generally used. These dual lumen catheters are usually configured so that there is a shorter lumen that aspirates blood from a blood vessel of a patient to a dialysis machine where it is processed for the removal of toxins, and a longer lumen that infuses the purified blood to the patient. The shorter lumen utilized for aspiration is generally referred to as the \u201carterial lumen,\u201d while the longer lumen utilized for infusion is generally referred to as the \u201cvenous lumen.\u201d The reason for the different lengths is to minimize co-mingling of aspirated and infused blood.\nThe primary problems occurring in dual lumen dialysis catheters are blood clotting (thrombosis) and fibrin (the protein formed during normal blood clotting that is the essence of the clot) sheath formation. Thrombus and fibrin sheath formation can occlude distal tips of the dialysis catheter lumens, resulting in loss of catheter function when such an occlusion prevents blood flow. This typically occurs initially in the arterial lumen used for aspiration of blood from a patient. A secondary problem relates to the arterial lumen \u201csucking\u201d against the vessel wall, in which the arterial lumen openings become fully occluded by the patient's vasculature.\nTo specifically address these problems, a new type of dialysis access catheter has been designed that utilizes independent \u201cfree floating\u201d distal tip sections that separate at a distal end of the catheter to theoretically reduce the likelihood of potential occlusion and \u201csucking\u201d during dialysis treatment. U.S. Pat. No. 6,001,079 to Pourchez and U.S. Pat. Nos. 5,947,953 and 6,190,349 to Ash et al., all incorporated by reference herein, are directed to said new type of catheter, hereinafter referred to as a \u201csplit-tip catheter.\u201d\nCatheters used for hemodialysis are generally inserted according to standard technique. This includes identifying the location of the central veins in which the catheter is to be inserted, utilizing the Seldinger technique to cannulate the central vein with a guidewire, and passing various instruments over the guidewire until the catheter is in place. In the case of a small diameter catheter made of relatively stiff material (i.e., a percutaneous catheter for acute access), the catheter itself can be passed over the guidewire and into the central vein. In the case of a larger diameter catheter, the passage to the central vein must be widened through the use of dilator(s), after which the catheter can be passed over the guidewire. Another possibility available to the physician is to use an introducer sheath, which is passed over the guidewire (after dilation or simultaneously therewith) forming a path for the catheter to access the central vein.\nDue to the configuration of the split-tip catheter, standard insertion techniques must be modified. With respect to non-sheath insertions, the separate distal tip sections must each be passed over the guidewire to ensure their arrival at the targeted central vein without attendant problems. Suggested means for accomplishing this (in a dual lumen catheter) is to pass the back end of the guide wire through the distal end hole of one tip section and out a side hole thereof into the distal end hole of the other tip section (Patel et al., J. Vasc. Interv. Radiol. (2001), vol. 12, pages 376-378). However, such a method can become problematic depending on the placement of the side holes and, in fact, may be impossible in the event that no side holes are provided. With respect to sheathed insertions, the separate distal tip sections may tend to flare outward and catch or \u201csnag\u201d on the inner wall of the delivery sheath thereby preventing smooth delivery of the catheter to the central vein.\nThere is a need for improvements to split-tip catheters and new configurations thereof, which enhance the design and further alleviate concerns with regard to typical problems encountered in hemodialysis catheters. There is also a need for new configurations and new methods for insertion of split-tip catheters to overcome problems associated therewith"}
-{"text": "This invention relates to an improved load cell of the type used to measure tension or compression loads and more particularly relates to a washer type load cell adapted to be used with gap-spanning sensor elements.\nLoad cells have many applications and are used in the prior art in electronic weighing systems for monitoring force and pressure and so on. As such, these cells employ a suitable sensor element such as piezoelectric, reluctance, capacitive or piezoresistive element which is positioned in a housing to monitor the force or pressure exerted on the housing when the load cell is positioned in a suitable environment.\nThe most popular types of load cells employ a strain gage element which may be an unbonded or bonded wire gage or a foil or a bonded semiconductor element. These gages provide higher sensitivity while further eliminating mechanical linkages and moving parts. They also enable one to achieve good temperature specifications and serve to reproduce load deflection in the elastic member or housing with good linearity.\nA major concern in the design of the load cell resides in the housing arrangement. In this manner, the housing serves to couple or transmit the force or pressure to be monitored by the cell to the sensor associated with the cell. The housing would therefore be capable of transmitting such forces to the sensor in a linear manner without providing undue distortions and improper bending movements when the monitored load is not properly aligned with the axis of the load cell. Some load cells are required to have high overall stiffness. Thus, the load cell housing should provide a high overall stiffness while still allowing sufficient strain to be monitored by the associated sensor element.\nIt is therefore an object of the present invention to provide a load cell transducer capable of accurately monitoring an external force employing a housing which is simple and relatively inexpensive to fabricate; and which provides a relatively stiff structure, but also allows for adequate strain levels to be transmitted to the sensor element or elements."}
-{"text": "A NAND type flash memory with a three-dimensional structure manufactured using a BiCS (Bit Cost Scalable Flash memory: BiCS) manufacturing technique is referred to in the art as BiCS memory."}
-{"text": "To meet the demand for wireless data traffic, which has increased since deployment of 4th-generation (4G) communication systems, efforts have been made to develop an improved 5th-generation (5G) or pre-5G communication system. Therefore, the 5G or pre-5G communication system is also called a \u2018beyond 4G network\u2019 or a \u2018post long-term evolution (LTE) system\u2019.\nIt is considered that the 5G communication system will be implemented in millimeter wave (mmWave) bands, e.g., 60 GHz bands, so as to accomplish higher data rates. To reduce propagation loss of radio waves and increase a transmission distance, a beam forming technique, a massive multiple-input multiple-output (MIMO) technique, a full dimensional MIMO (FD-MIMO) technique, an array antenna technique, an analog beam forming technique, and a large scale antenna technique are discussed in 5G communication systems.\nIn addition, in 5G communication systems, development for system network improvement is under way based on advanced small cells, cloud radio access networks (RANs), ultra-dense networks, a device-to-device (D2D) communication, a wireless backhaul, a moving network, a cooperative communication, coordinated multi-points (CoMP), a reception-end interference cancellation, and the like.\nIn the 5G system, a hybrid frequency shift keying (FSK) and quadrature amplitude modulation (QAM) modulation (FQAM) and a sliding window superposition coding (SWSC) as an advanced coding modulation (ACM) scheme, and a filter bank multi carrier (FBMC) scheme, a non-orthogonal multiple Access (NOMA) scheme, and a sparse code multiple access (SCMA) scheme as an advanced access technology have been developed.\nEnhancement of camera performance and implementation of a personal broadcast service application facilitate live broadcast with a high-definition image of definition equal to or greater than 2 k ultra high definition (UHD) on a client terminal such as a smart phone, and/or the like.\nFor transmitting a high-definition image in real time in a client terminal, a sufficient data rate, e.g., a data rate equal to or greater than 6 Mbps needs to be guaranteed in an uplink channel. However, a bandwidth of a wireless channel which is based on a cellular network (hereinafter, it will be referred to as \u201ccellular channel\u201d) is narrow. Thus, when the number of client terminals in a serving evolved node B (eNB) increases, the data rate per client terminal decreases thereby making it difficult to guarantee a high data rate.\nMeanwhile, a wireless communication system has proposed various technologies for increasing transmission capability, and a typical one is a carrier aggregation (CA) technology. The CA technology has been implemented with various forms, and an inter-heterogeneous eNB CA technology as one of the various forms is a technology in which frequency resources of different communication networks may be aggregated. For example, the inter-heterogeneous eNB CA technology may aggregate frequency resources of an LTE mobile communication system. Further, the inter-heterogeneous eNB CA technology may aggregate a frequency resource of an LTE mobile communication system and a frequency resource of a third generation (3G) mobile communication system, or a frequency resource of an LTE mobile communication system and a frequency resource of a wireless local area network (WLAN), e.g., a Wi-Fi, and/or the like. So, the inter-heterogeneous eNB CA technology may effectively use resources in various communication schemes for increasing a data rate.\nAnother one of the various technologies for increasing the transmission capability proposed in the wireless communication system is a multi-path transmission control protocol (MPTCP). The MPTCP is an example of a new transmission layer protocol which is generated for increasing a data rate using a multi-path. The MPTCP may combine a plurality of transmission paths connecting hosts, e.g., subflows to form one multi-path, e.g., an MPTCP flow and exchange data through the plurality of transmission paths. So, if the MPTCP is used, a preset number of transmission paths may be generated between the hosts.\nFor example, if the MPTCP is used, a data rate sufficient for providing a live video with high definition in real time may be acquired by combining insufficient cellular channel resources, e.g., uplink cellular channel resources. Further, hosts which transmit and receive data based on the MPTCP may distribute traffic to other transmission paths if a congestion situation occurs on a specific transmission path or a connection thereof is released due to various reasons.\nIn a general communication system supporting a multi-path, there is a need for a congestion control scheme for maximizing a data rate between hosts. However, in a current MPTCP standard, a congestion control scheme for an MPTCP is not defined, so a congestion control scheme for a transmission control protocol (TCP) is used. Due to this, congestion control in which a characteristic of an MPTCP is reflected is not performed, so various issues such as fairness violation due to duplicative occupation of a bottleneck in a multi-path, packet reordering due to difference among data rates and difference among round trip time (RTT) of transmission paths, and/or the like may occur.\nThe above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
-{"text": "1. Field of the Invention\nEmbodiments of the present invention relate to a coating that is applied to a surface. More particularly, embodiments of the present invention relate to a protective coating that is applied to the underwater portion of a marine vessel so as to inhibit the growth of marine foulants.\n2. Description of the Related Art\nMarine vessels that reside in a water environment over certain lengths of time can accumulate biological growth, known as foulants, on those surfaces that are in contact with the water. Diverse species of hard and soft fouling organisms, such as barnacles, zebra mussels, algae, and slime, form colonies on the underwater surfaces of the vessel, particularly when a vessel is docked, because each requires a permanent anchorage in order to mature and reproduce. Marine growth fouling adds weight to a ship, increases the amount of fuel consumed, and reduces its speed.\nHistorically, to combat the growth of marine foulants, the underwater surfaces of ships have been coated with antifoulant paints, which often include toxic materials to inhibit biological growth. The antifoulant paints may degrade and break down over time, releasing the toxic materials from the marine vessel into the surrounding water. These toxic materials may include volatile organic compounds (VOCs) and hazardous air pollutants (HAPs). The International Maritime Organization and the United States Environmental Protection Agency have enacted regulations and standards that restrict the emission of VOCs and HAPs from antifoulant paints. The decomposition and break down of the antifoulant paint results in reduced efficacy of the protection afforded by the antifoulants, thereby requiring reapplication of the paint in a relatively short time. Thus, a coating material is required that can be applied to the underwater surfaces of a marine vessel which repels the growth of fouling organisms on such surfaces and has an extended lifetime without releasing significant amounts of toxic materials into the environment."}
-{"text": "Typically, most of the power interruptions to which lamps are subjected are of very brief duration, e.g., under one second, but even an interruption of this brief duration will cause a typical HID lamp to be extinguished. If the usual HID lamp is subjected to such a brief power interruption, it does not immediately reignite upon restoration of normal voltage thereacross. Typically, some cooling of the gases within the HID lamp will be required before the lamp will reignite in response to voltage applied thereacross. When the HID lamp does reignite in response to such voltage, its lumen output at the time of reignition is only a small fraction of normal. Gradually, this lumen output increases until the lamp regains its full brightness."}
-{"text": "The IEEE standard, \u201cIEEE Std 1394-1995 Standard For A High Performance Serial Bus,\u201d is an international standard for implementing an inexpensive high-speed serial bus architecture which supports both asynchronous and isochronous format data transfers. Isochronous data transfers are real-time transfers which take place such that the time intervals between significant instances have the same duration at both the transmitting and receiving applications. Each packet of data transferred isochronously is transferred in its own time period. An example of an ideal application for the transfer of data isochronously would be from a video recorder to a television set. The video recorder records images and sounds and saves the data in discrete chunks or packets. The video recorder then transfers each packet, representing the image and sound recorded over a limited time period, during that time period, for display by the television set. The IEEE 1394-1995 serial bus architecture provides multiple channels for isochronous data transfer between applications. A six bit channel number is broadcast with the data to ensure reception by the appropriate application. This allows multiple applications to simultaneously transmit isochronous data across the bus structure. Asynchronous transfers are traditional data transfer operations which take place as soon as possible and transfer an amount of data from a source to a destination.\nThe IEEE 1394-1995 standard provides a high-speed serial bus for interconnecting digital devices thereby providing a universal I/O connection. The IEEE 1394-1995 standard defines a digital interface for the applications thereby eliminating the need for an application to convert digital data to analog data before it is transmitted across the bus. Correspondingly, a receiving application will receive digital data from the bus, not analog data, and will therefore not be required to convert analog data to digital data. The cable required by the IEEE 1394-1995 standard is very thin in size compared to other bulkier cables used to connect such devices. Devices can be added and removed from an IEEE 1394-1995 bus while the bus is active. If a device is so added or removed the bus will then automatically reconfigure itself for transmitting data between the then existing nodes. A node is considered a logical entity with a unique address on the bus structure. Each node provides an identification ROM, a standardized set of control registers and its own address space. Because of these advantages, the IEEE 1394-1995 standard provides for a unique networking structure that is capable of incorporating audio/video devices, media play/record devices and computing/display devices.\nA diverse range of products can be implemented with the ability to connect to an IEEE 1394-1995 serial bus network. These devices can have capabilities and functionality ranging from very simple to very complex. Specifically, a variety of audio/video devices, media play/record devices and computing/display devices are capable of being linked together over an IEEE 1394-1995 serial bus networking structure to support asynchronous and isochronous data transfers between the devices.\nThe IEEE 1394-1995 serial bus allows a collection of devices to work together in a high bandwidth, distributed environment to maximize the overall efficiency and functionality of the network. This allows manufacturers to remove expensive pieces of functionality from one device and locate that functionality in another device on the network, instead of duplicating this functionality in all devices on the network. While some of the devices have limited functionality and are relatively inexpensive, such devices require the support and interaction of other devices in order to bring the full functionality of the devices within the network to the user.\nThe AV/C Digital Interface Command Set is a command set used for data transactions between consumer audio/video equipment over an IEEE 1394-1995 serial bus. Neither the IEEE 1394-1995 serial bus nor the AV/C command set provide a master-slave relationship between the devices coupled within the IEEE 1394-1995 serial bus network. Instead, both the IEEE 1394-1995 serial bus and the AV/C command set operate based on a cooperative peer-to-peer coexistence of devices within the network. Discrete AV/C command and response data packets are transferred between networked devices over an IEEE 1394-1995 serial bus in an asynchronous data stream. The AV/C command and response data packets are typically formatted according to the AV/C protocol outlined in the AV/C Digital Interface Command Set. Transfers of AV/C command and response data packets over the IEEE 1394-1995 serial bus network are supported by an AV/C architecture. The AV/C architecture is comprised of lists and tables that help devices create, process and/or transmit AV/C command and response data packets. The AV/C architecture includes an AV/C bulletin board subunit that is typically dedicated to a resource device subunit such as a tuner, receiver or recoding unit.\nAn AV/C bulletin board subunit is an information structure that is shared between devices networked over an IEEE 1394-1995 serial bus network. A resource schedule bulletin board is also an information structure that supports information shared between coupled devices within a network. The resource schedule bulletin board provides the organizational structure through which shared data is organized and communicated. The resource schedule bulletin board contains sub-boards of lists with entry descriptors that represent encoded data that can be shared between devices within the network via descriptor commands. A dedicated AV/C bulletin board subunit typically supports the information architecture between that device and all compatible posting devices within an IEEE 1394-1995 serial bus network. A posting device writes a request entry to a write enabled list within the resource schedule bulletin board specifying when it will use the resource.\nNeither the IEEE 1394-1995 serial bus nor the AV/C Command Set provide a master-slave relationship between the devices coupled within the IEEE 1394-1995 serial bus network. Instead, both the IEEE 1394-1995 serial bus and the AV/C Command Set operate based on the cooperative peer-to-peer coexistence of devices within the network transmitting data formatted in accordance with the AV/C protocol. The communicating devices exchange command and response data directly with each other without the intervention of a systems resource manager. Each device or device subunit is responsible for managing scheduling affairs and storing resource requests.\nWithin the current AV/C protocol, resource time of a resource device is made available to any requesting control device. Within this protocol a device competing for use of a resource device can acquire the right to use the resource device by reserving the resource device, even though another client device may have previously submitted an entry to the resource schedule board requesting the resource device at the same time. Devices may modify or delete data contained within the schedule entries submitted by other requesting control devices without acknowledging to the other requesting control devices that their respective schedule entries are being or have been modified. Thus, within the current AV/C protocol there is no method to allow control devices within a network to monitor schedule data or to provide a peer-to-peer review of schedule modifications."}
-{"text": "1. Field of the Invention\nThe present invention relates to a rotor employed in a flywheel energy storage device. More specifically, the present invention relates to a rotor wherein the power density and the energy density of the system are optimized while the manufacturing cost is minimized. Methods for manufacturing a rotor according to the present invention are also disclosed.\n2. Description of the Prior Art\nFlywheel systems currently being designed for mobile energy storage are generally intended to replace batteries in electrically powered vehicles. One specific example of such a system was disclosed in commonly assigned, co-pending application Ser. No. 08/148,361, which is incorporated herein by reference for all purposes. In such applications, multiple units are needed to store the required energy, so that each motor-generator need supply only a small portion of the vehicle's power. In systems where all of the surge power must be supplied by a single flywheel, the relatively large size of the single motor-generator makes it difficult to provide the needed energy density without reducing safety factors, e.g., for radial stresses, to unacceptable low levels or raising manufacturing costs to exorbitantly high levels.\nU.S. Pat. No. 3,741,034 discloses rotor designs using high strength-to-weight ratio filament wound composites in relatively thin concentric cylinders, which cylinders are separated by radial springs. While this arrangement limits the radial stresses to tolerable values, it is expensive to manufacture.\nU.S. Pat. No. 3,859,868 discloses techniques for varying the elasticity-density ratio of the rotor elements to minimize radial stresses. On the other hand, U.S. Pat. Nos. 4,341,001 and 4,821,599 describe the use of curved metallic hubs to connect the energy storage elements to the axle. Additionally, U.S. Pat. No. 5,124,605 discloses a flywheel system employing counter-rotating flywheels, each of which includes a hub, a rim and a plurality of tubular assemblies disposed parallel to the hub axis for connecting the hub to the rim while allowing for differential radial expansion between the hub and the rim.\nNone of these references deal with the integration of a large, high power motor-generator into the flywheel energy storage system currently being designed for vehicles."}
-{"text": "Structured wiring enclosures or \u201ccans\u201d are typically recessed in walls. Known enclosures have limited capability and versatility with regard to mounting electronic devices and/or cable or wire management.\nEnclosure door covers or enclosure doors typically mount to an underlying enclosure. Known covers or doors for structured wiring enclosures do not provide substantial functionality beyond covering and providing access to the enclosure."}
-{"text": "1. Field of the Invention\nThis invention relates generally to actuators for ergonomic systems and, more particularly, to actuator systems for seat adjustments.\n2. Related Art\nErgonomic supports for seats, such as lumbar and bolster systems, are typically moved by means of actuators that can be operated by hand or driven by a motor. Four-way power lumbar devices, with an arching mode and a translation mode, have traditionally required a motor for each mode of operation. In comparison, the invention set forth by U.S. Pat. No. 6,050,641 is a four-way power lumbar system that requires only a single-motor and reduces the duplication of gearbox components. Prior to this invention, U.S. Pat. Nos. 5,197,780 and 5,217,278 had only described four-way lumbar devices that were manually operated by a single control knob.\nWhile these devices are an improvement over the conventional four-way lumbar devices that required multiple motors or multiple control knobs, the highly competitive markets for furniture and automotive seats place a premium on continued optimization of devices that provide comfort and convenience for seat occupants. In particular, there is a need for improved actuation systems that are less prone to failures and more efficiently transfer power to the actuators. For example, with regard to the manually-operated four-way lumbar systems, the gearing systems are inefficient because, in addition to the gears, they have a transmission that requires at least one non-gear alignment device to maintain the proper engagement between the driver gear and the driven gears. Therefore, such devices have an alignment device that is necessary for the convenient operation of the lumbar system and increase the potential for a failure in the system.\nAdditionally, there is a need for actuation systems that are more modular, increasing the commonality of parts between two-way and four-way power lumbar devices, manually-operated and motor-driven actuation systems, and lumbar systems and bolster systems. For example, the prior art actuation units and driven gears are designed to fit within a single housing along with the driver gears and the transmission system, thereby limiting the range of motion that is capable for the actuation units themselves. Different types of lumbar devices and bolster devices are typically designed to provide different levels of support and often require different levels of actuation, thereby affecting the size of the actuators. The prior art devices do not easily allow for changing the actuators according to various sizing requirements, and the confined housing could prevent the same actuation system from being used for different lumbar devices or for a lumbar device and a bolster device. Therefore, entirely different actuation systems would need to be designed, and separately manufactured, depending on the actuators' usage, range of motion, and sizing."}
-{"text": "A technology has been conventionally known that reduces the amount of energy consumption of a service side by altering a service timetable so as not to exceed a certain amount of electric power that is under contract with an electric power company. In such a technology, for the purpose of energy saving of a whole of a railroad system including station facilities, the energy saving of a station side may also be performed, in addition to the energy saving on the service side. In this case, it is common that the energy saving on the service side and the energy saving on the station side are performed independently from each other.\nIn such a technology described above, for an example, it is desirable that the energy saving on the service side and the energy saving on the station side be able to be performed in a unified manner."}
-{"text": "Websites often provide a social sharing feature in the form of a share widget for content hosted by the website. These share widgets enable users to quickly share content with members of their social network using multiple social network services (e.g., Facebook\u00ae, Twitter\u00ae, Pinterest\u00ae, or the like). Although such functionality may provide users with a link between the website's services and the social network services, the websites themselves are often unable to ascertain which of their users correspond to the social network profiles used to publish content from the websites."}
-{"text": "In a network device, a scheduler is commonly used to schedule the processing of data units (e.g. packets, cells, frames, etc.) associated with multiple traffic flows. Typically, a weighted round-robin (WRR) scheduler may process data units according to WRR weights associated with queues. Each WRR weight may indicate a level of priority, in which the larger the WRR weight, the higher the priority. For example, consider three queues, namely, queue 0, queue 1, and queue 2, which are assigned respective WRR weights of 10, 20, and 30. Given this configuration, the WRR scheduler services queue 0 for 10 timeslots, and queue 1 may have to wait those 10 timeslots before the WRR scheduler services queue 1. Additionally, queue 2 may have to wait 30 timeslots (i.e., 10 timeslots for queue 0 and 20 timeslots for queue 1) before the WRR scheduler services queue 2. As a result, the WRR scheduler may cause periods of burstiness to the traffic flow associated with, for example, queue 2, which may increase the amount of buffering and/or may cause a higher latency to occur with respect to queue 0 and/or queue 1."}
-{"text": "1. Field of the Invention\nThe present general inventive concept relates to a computer system and a program restoring method thereof, and more particularly to a computer system storing with a restoration program and a restoring method using the restoration program.\n2. Description of the Related Art\nIn general, a computer system performs various functions through hardware such as a central processing unit (CPU), and software such as an operating system, a device driver, and an application. In the case of the software, update is often periodically provided by an external server or the like with development of its technology. However, the updated software may have a problem with its operation due to various factors, for example, if it is infected with a virus while using a computer system, deleted by mistake, modified, or etc. At this time, the updated software having the problem has to be restored or recovered, but a restoration solution initially provided by a manufacturer of a computer system provides only a previous version of software. Thus, a user has to access the external server or the like after restoring the software, and then update the restored software with the latest version."}
-{"text": "The present disclosure relates to providing and enforcing security procedures in computer systems in general and more specifically relates to providing and enforcing security procedures in mobile information systems. In a mobile information system, subscribers may register with a wireless service provider to receive various types of content from the service provider on their mobile devices. The subscriber's mobile device may include a resident interactive multimedia application environment that includes capabilities for displaying graphics, video, animation, audio, and the like. Examples of such interactive multimedia application environments are the different versions of the Adobe\u00ae Flash\u00ae-based platform. Content provided to mobile devices equipped with such application environments is sometimes delivered in executable file formats such as the precisely described SWF (small web format) binary vector graphics format. SWF provides a compact, TAG-based, easily extendible format that supports streaming, bitmap and vector graphics, and scripting."}
-{"text": "Many websites and applications apply various techniques and methods to discern a human user from automated software prior to granting certain access or privileges to the user. Distinction between a human user and automated software is often desired, such as when a user is signing up for a service, when a user is attempting to log in to a service or account, when a user is participating in an online poll or forum, or when a user generating an online post. Many websites employ CAPTCHAs (Completely Automated Public Turing tests to tell Computers and Humans Apart) to achieve such functionality. A common CAPTCHA includes a depiction of distorted text images and requires transcription of the distorted text image. However, websites, applications, and other services that implement such CAPTCHAs risk inadvertently driving away traffic by legitimate users (i.e. non-spamming humans) while providing only minimal protection against automated software and human spammers. Alternatively, websites, applications, and other services that fail to implement some such security feature risk security breaches.\nTherefore, there is a need in the human verification field to create a new and useful method for generating a human likeness score indicative of the legitimacy of a use. This invention provides such a new and useful method."}
-{"text": "Press-hardened steel alloys are being used for sheet-metal parts incorporated in vehicle body structures that may be assembled together with rivets or welding. One example of a press-hardened steel is boron steel sold under the designation Usibor\u00ae 22MnB5. Press-hardened steel can be water cooled or oil cooled to a desired level of hardness from 450 to 520 HV. Press-hardened steel may be annealed to reduce the hardness to 140 HV.\nPress-hardened steel parts may be assembled to other steel parts by welding. But, new automotive assemblies may include combinations of parts made of different materials such as aluminum and composite parts. An Ultra High Strength Steel (UHSS) beam formed by press hardening and a composite part or an aluminum part cannot be efficiently joined together in a welding operation. The preferred technique for joining such part assemblies is to rivet or otherwise fasten the parts together. The hardness of such high strength parts poses significant challenges in high volume manufacturing operations because the rivets have difficulty penetrating the UHSS beam."}
-{"text": "This invention relates to a toy vehicle and, more particularly, to a self-propelled toy vehicle in the shape of a robot which is propelled while performing predetermined tasks.\nIn the field of robot-like toy vehicles, there is an endless demand for novel performance. The present invention is directed to satisfying this demand."}
-{"text": "International Publication No. WO2010/146738 discloses a shift register for scanning gate bus lines. This shift register is composed of a plurality of cascade-connected stages the number of which corresponds to the number of the gate bus lines. Each stage has terminals to which a high-level power source voltage and a low-level power source voltage are input, respectively, and terminals to which an output signal sent from the previous stage and an output signal sent from the next stage are input, respectively. Further, each stage has a terminal to which a clock signal is input, and a terminal through which an output signal of the current stage itself is output as a gate pulse. Each stage is composed of a plurality of transistors including a transistor (M5) and a transistor (M6), and a capacitor that are connected to one another. The transistor (M5) is connected to the terminal to which the high-level power source voltage is input. Further, the transistor (M6) is connected to a terminal into which the low-level power source voltage is input.\nAn internal node in a shift register in each stage has a potential that changes with a clock signal. According to a change of the potential of the internal node, the transistor (M5) is turned on, and a power source voltage VDD is output as a gate pulse via the transistor (M5). This shifts a gate bus line into a selected state. Further, the transistor (M6) is turned ON in response to the clock signal, and the transistor (M6) causes the potential of the internal node to have a power source voltage VSS, thereby decreasing the gate pulse. This shifts the gate bus into a non-selected state."}
-{"text": "Triple negative breast cancer (TNBC) represents a significant proportion (about 20-25%) of breast cancer patients. TNBC is characterized by the absence of HER2, estrogen receptor (ER), and progesterone receptor (PR). TNBC has a poor prognosis, and no targeted approach to a therapy has been found to date."}
-{"text": "The invention relates to a method of controlling the torque developed at a belt shaft of a belt retractor coupled to an electric motor. DE 100 18 972 discloses a belt retractor whose belt shaft is driven by an electric motor rather than by a coil spring. The electric motor is coupled to the belt shaft by a toothed belt. Assuming the electric motor to be supplied with a constant amount of current, friction and the flexing work of the toothed belt result in a hysteresis of the seat belt forces, being different when the seat belt is wound up and when it is unwound.\nIt is an object of the invention to essentially avoid hysteresis forces in order to keep the pressure exerted by the seat belt on the body of the vehicle occupant essentially constant.\nIn the method according to the invention, torque developed at a belt shaft of a belt retractor coupled to an electric motor is controlled by detecting a direction of rotation of the belt shaft and, depending on the direction of rotation as detected, changing electric current supplied to the electric motor from a first value to a second value that is different from the first value by an amount required to compensate for hysteresis.\nAccordingly, to equalize the seat belt forces the amount of electric current supplied to the electric motor is adjusted in such a way that it is increased in the winding direction and decreased in the unwinding direction, resulting in substantially equal torque at the belt shaft and hence essentially equal seat belt forces when the seat belt is wound up and unwound.\nIn accordance with a preferred embodiment of the invention, a standstill of the belt shaft is also detected. The seat belt forces needed to keep the pull-out length of the seat belt constant are smaller than the forces needed to unwind it. Accordingly, to compensate for the smaller forces, a higher amount of electric current is supplied to the electric motor.\nIn case of an emergency, when strong deceleration forces act on the vehicle, inertia forces tend to move the body of the vehicle occupant away from the seat, entraining and unwinding thereby the seat belt at a speed that is higher than during normal movement of the vehicle occupant. In order to protect the vehicle occupant, such a rapid unwinding of the seat belt has to be avoided.\nIn accordance with another preferred embodiment of the invention, rotational speed of the belt shaft is measured in an unwinding direction and when a predetermined first threshold value of rotational speed is exceeded, the resistance against pulling out the seat belt is increased by increasing the amount of electric current supplied to the electric motor, the rate of increase being constant. Only when a second threshold value of rotational speed, corresponding to a threshold value of the amount of electric current, is reached, a mechanical locking mechanism on the belt retractor is activated. Otherwise a restriction of the seat belt is achieved by means of the electric motor only.\nA preferred use of the inventive method is in a retractor of the type comprising a frame, a belt shaft rotatably fixed in the frame and an electric motor coupled to the belt shaft by a drive belt, wherein the direction of rotation can be detected and, depending on the direction of rotation as detected, the motor current is adjusted."}
-{"text": "Poly(para-phenylene) (PPP) is a fully aromatic, rigid rod polymer with unique structural and conductive properties. As an engineering plastic, its attractiveness arises from its thermal stability (mp&gt;500.degree. C.), high strength, chemical inertness, and solvent resistance. When doped with either n- or p-type dopants, the polymer forms highly conducting charge transfer complexes with conductivities up to 500 S/cm. However, the structural properties which make PPP so attractive also make it a difficult polymer to synthesize. In addition, many of the observed properties of the polymer depend on the method of production.\nPrevious methods of producing PPP directly have met with only limited success. For example, oxidative cationic polymerization of benzene to produce PPP has been attempted. However, only short oligomers of ten to fifteen repeat units containing mixtures of linear 1,4- and non-linear 1,2-units were formed.\nPolymerizations using nickel catalyzed aryl coupling of 1,4-dihalobenzenes were attempted. While this method produced a completely linear molecule, only short oligomers consisting of ten to twelve units were formed.\nThe problem with these direct synthetic methods is that the inherent insolubility of the polymer causes it to precipitate out of solution before high molecular weight materials can be formed. Electrochemical coupling of benzene has also been used, but the resulting film is insoluble and composed of a mixture of 1,4- and 1,2- units.\nIn order to circumvent the problem of the inherent insolubility of PPP in production and processing, soluble precursor methodologies have been developed. For example, polymers of 1,3-cyclohexadiene (CHD) have been used as a soluble precursor polymer. In particular, poly(cyclohexadiene) has been reacted with bromine and then pyrolyzed to eliminate HBr. Unfortunately, this polymerization route also produces a precursor polymer with a mixture of 1,4-and 1,2-linkages. In addition, the elimination reaction is not very efficient, since HBr readily reacts with unsaturated intermediates.\nRecently, the efficient production of PPP has been reported, via the pyrolysis of a soluble precursor polymer prepared from the radical polymerization of the acetyl and methoxycarbonyl derivatives of 5,6-dihydroxy-1,3-cyclohexadiene (DHCD); see, e.g., D. G. Ballard et al., Macromolecules, Vol. 21, pp. 294-304 (1988) and D. R. McKean, Macromolecules, Vol. 20, pp. 1787-1792 (1987). The starting cis-diol is produced by the microbial oxidation of benzene. The precursor films are soluble and can be processed before pyrolyzing into the final polymer. However, the radical polymerization produces about 85% 1,4-units and 15% 1,2-units. The 1,2-units create \"kinks\" in the polymer, thereby reducing the elimination efficiency of the precursor and the mechanical properties of the final polymer.\nIn order for the good mechanical properties of PPP to be realized, an aspect ratio of at least 100 consecutive linear units per 1,2-unit must be obtained. All previous routes to PPP have either produced low molecular weight materials due to insolubility of the growing polymer, or have incorporated a significant amount of 1,2-linkages in the chains, or both. Hence, it is desirable to find an exclusively 1,4-polymerization method which can be used in combination with the efficient precursor method described above."}
-{"text": "The present invention relates to a process for the extraction of cesium ions from aqueous solutions using an adduct compound comprising a macrocyclic polyether (crown ether) and an inorganic heteropoly acid component.\nThe processing and solidification of medium active aqueous wastes (MAW), developed during the reprocessing of irradiated nuclear fuels and/or nuclear breeder materials, e.g., in form of waste solutions, waste concentrates or waste sludges, is rendered more difficult by the presence of small quantities of the cesium isotopes Cs.sup.134 and Cs.sup.137. During the vitrification of MAW, cesium evaporates to a noticeable degree and is also easily leached out of the solidified products intended for permanent storage, such as glass blocks, cement blocks and bitumen blocks. Selective extraction of the Cs would considerably simplify the treatment of medium active wastes. Furthermore, if Sr.sup.90 were simultaneously extracted from the MAW, low radioactive aqueous wastes (LAW=low active waste) would be obtained since, after a relatively short decay period, practically the entire activity of MAW originates from the relatively long-lived radionuclides Cs.sup.137 (t.sub.1/2 =30 yrs) and Sr.sup.90 (t.sub.1/2 =26 yrs).\nAccording to prior-art processes, the extraction of Cs and Sr was achieved mainly by means of coprecipitation reactions. The coprecipitation reactions, however, did not result in satisfactory DF (decontamination factor) values for Sr and Cs. Therefore, different processes have been investigated which would make possible a selective extraction of the two radionuclides, either jointly or individually.\nThe extraction processes for Cs.sup.+ ions developed thus far, however, are not suited for the extraction of Cs.sup.+ from a typical MAW with a high concentration of NaNO.sub.3 and free nitric acid. The composition of such a typical MAW is shown in the following Table 1:\nTABLE 1 ______________________________________ Components Concentration (strongly acidic MAW) g/dm.sup.3 ______________________________________ Nitric acid 64 (1 mol/dm.sup.3) Sodium nitrate 42.5 (0.5 mol/dm.sup.3) Uranium 1.0 Lead 0.8 Calcium 0.2 Magnesium 0.2 Iron 0.15 Ruthenium 0.07 Potassium 0.05 Zinc 0.05 Copper 0.04 Manganese 0.02 Chromium 0.02 Nickel 0.01 Zirconium 0.005 Cesium 0.004 Cerium 0.004 Strontium 0.002 Antimony 0.002 Niobium 0.001 ______________________________________\nThe extraction processes for Cs.sup.+ ions that have been developed use mostly organic reagents which are added either to the organic phase or to the aqueous phase to be extracted. These reagents, however, would be destroyed either hydrolytically or by oxidation, or would be made ineffective by protonization if they came into contact with the MAW whose essential components are such as those described in Table 1 above.\nCs.sup.+ selective heteropoly acids are known which are very stable in the highly acid environment provided by an MAW having a composition such as described in Table 1. The Cs.sup.+ selective heteropoly acids are used as extraction agents in the form of an organic phase with nitrobenzene as a solvent. For example, 12-molybdophosphoric acid (HPMo) is a Cs.sup.+ selective heteropoly acid which can be dissolved in nitrobenzene. In this form, the heteropoly acids are equally unsuitable for the continuous Cs.sup.+ extraction from an MAW of the type described in the above Table 1, since they bleed considerably from the nitrobenzene into the aqueous phase due to their good water solubility. They would have to be replaced continuously.\nIf HPMo is used, the molybdenum would enter into the MAW in large quantities and would furthermore interfere with the subsequent processing of the MAW, such as during vitrification. Finally, depending on the Cs.sup.+ concentration of the MAW and the heteropoly acid concentration in the nitrobenzene phase, there is the danger of precipitate formation which could lead to considerable disruptions in a continuously operating system.\nIn contrast to the above prior-art water soluble extraction agents for Cs.sup.+, dibenzo crown ethers have a very low water solubility and are thus, at least in part, suitable for use in a continuous extraction process.\nSince crown ethers are neutral ligands, the anion in the aqueous solution being treated is also extracted during the extraction. This, however, significantly influences the extraction coefficient. In this way, simple cesium salts, such as chloride and nitrate salts, are very difficult to extract with the use of crown ethers in polar solvents, with the exception of nitromethane solutions and nitrobenzene solutions. For this reason, the prior art started to combine crown ethers with known Cs.sup.+ specific reagents which contain voluminous, polarizable anions. Up to now, this has been done by adding the Cs.sup.+ specific reagents, in most cases picric acid to the aqueous phase to be extracted,, and extracting with the aid of crown ethers, in most cases dibenzo-18-crown-6 (DB-18-C-6) or dicyclo-hexyl-18-C-6, in polar solvents, in most cases nitrobenzene.\nFor the extraction of Cs.sup.+ from nuclide containing solutions, J. Rais and P. Selucky proposed two extraction systems which use DB-18-C-6. In one method, as described in Czechoslovakian Pat. No. 149403, DB-18-C-6 was employed in an organic phase and dipicrylamine was added to the aqueous phase to be extracted. In a second method, as described in Czechoslovakian Pat. No. 149404, DB-18-C-6 was employed in an organic phase and sodium tetraphenylborate was added to the aqueous phase to be extracted.\nThese processes, however, are limited to alkaline Cs.sup.+ solutions having a pH of 11 to 13 inasmuch as dipicrylamine is protonized in the acid range, and sodium tetraphenylborate is hydrolyzed in the acid range. Moreover, both processes only function well in the absence of large quantities of Na.sup.+ and K.sup.+.\nFurthermore, the same authors, as disclosed in Czechoslovakian Pat. No. 165751, have produced an adduct of DB-18-C-6 with 12-molybdophosphoric acid which permits a column chromatographic extraction of the Cs.sup.+ from waste solutions. Large quantities of Na.sup.+ and K.sup.+, however, also present a problem here. The adduct of DB-18-C-6 with 12-molybdophosphoric acid was prepared by a two phase reaction in which an organic phase comprising a solution of the crown ether in dichloromethane was shaken with an aqueous phase comprising a solution of HPMo in 2 mol/1 HNO.sub.3. This patent also refers generally to adducts with 12-tungstophosphoric acid (HPW). A similar disclosure of an adduct of DB-18-C-6 with molybdophosphoric acid appears in J. Radioanal. Chem. Vol. 35 (1977), pages 351 to 359."}
-{"text": "A typical production-packaging apparatus for wrapping a stack of objects with a film has at least one film supply for the film, a film-feed head for feeding the film, and a device for pulling the film down over the stack of objects. The invention further relates to a method of wrapping a stack of objects with a film. Within the scope of the invention, \u201cfilm\u201d means in particular a plastic film, preferably an elastic thermoplastic film. It is within the scope of the invention that the stack of objects to be wrapped is a stack of products. The stack may also be a single object or whiteware, for example refrigerators or the like. The stack of objects to be wrapped is advantageously accommodated on a pallet.\nIn practice, it is not uncommon for stacks of objects or pallets carrying these stacks of objects having different orientations to be consecutively supplied to a wrapping station of an apparatus for wrapping the stack of objects. Thus, a stack of objects that in top view are not square or a pallet that in top view is not square having the length/width dimensions L/W is initially supplied and wrapped, and subsequently another stack of objects that has the dimensions W/L and that is rotated by 90\u00b0 about a vertical axis with respect to the first stack of objects must be supplied and wrapped. Additional measures are necessary here in order to obtain a film or film sheet that in such a case is also optimally adapted to the dimensions of the stack of objects or of the pallet.\nIn practice it is known, among other things, to lift the stack of objects using the elevator of a shrink-wrapping device and to rotate the stack of objects by 90\u00b0 for the covering operation. However, this often results the stack of objects falling over, together with product losses and facility down times. It is also known to provide in the conveyor a rotatable table underneath the film-feed head so that the stack of objects may be appropriately rotated. However, constrictions result in the transition regions between the rotatable table and the conveyor edge due to the fact that the pallets associated with the stacks of objects often become stuck at the transitions, and the stacks of objects then frequently fall over. It is also known in practice to design the entire film-feed head to be rotatable by 90\u00b0 so that the film can always be fed in the desired direction. This approach is relatively complicated, and the process is also too slow, in particular because of the large mass that has to be moved.\nIt is further known in practice for the film initially to be completely provided in the feed head and welded, and then transferred to a roller and rolled up by the roller. This roller is then rotated by 90\u00b0 with respect to the direction of transport, and the film is transferred from the roller to the pull-down device. These measures known in practice are also relatively complicated. The multiple transfers of the film result in a relatively high susceptibility to malfunctions."}
-{"text": "Embodiments generally relate to web based security management. More particularly, embodiments relate to the intelligent aggregation of platform device interaction information associated with web applications.\nEmerging markup languages such as HTML5 (Hypertext Markup Language 5, e.g., HTML5 Editor's Draft 8 May 2012, W3C), LLVM (e.g., LLVM 3.1, May 22, 2012, llvm*org), and other runtime or just in time (JIT) environment languages may support more robust multimedia related web platform development. The use of these advanced languages by an application developer, however, may also expose end user platform components such as graphics processors, memory, sensors, and so forth, to web applications, wherein the exposure of such components may lead to security concerns."}
-{"text": "The present disclosure relates to a projector. More specifically, the present disclosure relates to a presentation projector that dynamically masks out one or more objects positioned between the projector and a screen.\nDuring a typical office presentation, a digital data file such as a slideshow may be loaded onto a computing device and displayed on a screen via a projector. FIG. 1 illustrates an exemplary projector 100 projecting a presentation onto a screen 102. The projector 100 receives a video-in signal 104. The video-in signal 104 is processed by a liquid crystal display (LCD) panel controller 106. In a typical LCD projector, three LCD panels 108 are used, one each for projecting the colors red, blue and green. By opening or closing each pixel on the LCD panels 108, a variety of colors may be projected for each pixel by combining the outputs of the various LCD panels 108. Based upon the video-in signal 104, the LCD panel controller 106 determines the state of each pixel (i.e., opened or closed) for every pixel in the LCD panels 108. The LCD panel controller 106 transmits a control signal to the LCD panels 108 indicating the state (i.e., open or closed) for each pixel in the panel. A lamp 110 projects a strong beam of light through each of the LCD panels 108. The lamp 110 may be a high intensity lamp such as a metal halide lamp or other similar lamp. After the light emitted by the lamp 110 passes through the LCD panels 108, the light passes through one or more lenses 112, and the projected light 114 is directed by the lenses toward the screen 102.\nThe projector 100 may be a portable device placed on a desk or table at a low enough level that the light emitted from the projector may be projected at eye level of a person 116 walking in front of the projector. The person 116 walking in front of the projector may experience discomfort due to bright light shining in their eyes. Audience members may be distracted by the images projected onto the person 116 walking in front of the projector, as well as the shadow 117 cast by the person on the screen 102."}
-{"text": "Cycle exercise machines, sometimes also known as \"stationary bicycles\", have been known for some time. These exercise machines typically include a seat much like a bicycle seat, for supporting a person sitting upright on the machine, and a foot-powered crank arrangement mounted below the seat. This crank arrangement in prior-art cycle exercisers includes a pair of crank arms extending outwardly from a rotatable shaft, with foot pedals mounted at the free ends of the crank arms. The crank arrangement is mounted for receiving the feet of a person seated on the exercise machine. A drive chain or belt transfers the rotary motion of the crank arrangement to an energy-absorbing device such as a flywheel or rotor, which provides a load force against which the exerciser expends energy while pedaling the crank mechanism. One such cycle exerciser is disclosed in U.S. Pat. No. 4,188,030 to Hooper.\nBecause cycle exercisers of the prior art use a crank mechanism which the user pedals while doing work on the machine, these exercise machines cannot provide a constant transfer of power from the user to the load throughout the entire stroke of each pedal. Crank pedal mechanisms have a null at the top and bottom of each stroke, where the effective length of the pedal crank arms diminishes to zero length as the crank mechanism passes through the dead center position. The effective lever arm of the pedals then increases sinusoidally to the maximum effective length as the pedal crank arms pass through the 90.degree. position, i.e., half-way between the top to the bottom of each stroke. Because this effective lever arm constantly changes, the rate at which the exerciser can effectively expend energy doing work on the load device likewise changes throughout each stroke of the pedal. This variation lessens the possible maximum efficiency of cycle exercisers, as the rate of energy transfer for each stroke is maximized only momentarily during each stroke.\nAnother disadvantage of existing cycle exercisers, and in particular the crank mechanism used with those exercisers, pertains to the length of stroke. The stroke length is determined by the length of the crank arms used in the crank mechanism, and this length is not readily changeable. Changing the effective stroke length requires either replacing the entire crank assembly, or providing crank arms having variable positions for attaching the pedals on the crank arms. In either case, these alternatives are expensive and require the service of a technician to modify the length of the crank arms for users having significantly different lengths of stroke, namely, shorter or longer legs, or those of substantially different athletic ability."}
-{"text": "(1) Definition of Multi-Input Multi-Output (MIMO) Technology\nMIMO technology means a method for improving efficiency in data transmission and reception by using multiple transmitting antennas and multiple receiving antennas instead of a single transmitting antenna and a single receiving antenna. Namely, the MIMO technology is to increase capacity or improve throughput by using multiple antennas in a transmitter or receiver of a wireless communication system. Herein, multiple antennas will be referred to as MIMO.\nThe MIMO technology is an application version of a technology that a single message is completed by receiving data pieces from several antennas without depending on a single antenna path. The MIMO technology can improve a data transmission rate in a specific range or increase a system range for a specific data transmission rate. In this respect, the MIMO technology is a next generation mobile communication technology that can widely be used for a mobile station and a relay station. Also, the MIMO technology has received much attention as a next generation technology that can increase a transmission rate of mobile communication, which has reached the uppermost limit due to extension of data communication.\n(2) System Modeling in MIMO\nFIG. 1 is a schematic view illustrating a general MIMO communication system. As illustrated in FIG. 1, if the number of transmitting antennas is N7 and at the same time the number of receiving antennas increases to NR, channel transmission capacity increases theoretically in proportion to the number of antennas unlike that a plurality of antennas are used in only a transmitter or receiver. Accordingly, it is possible to improve a transmission rate and remarkably improve frequency efficiency. A transmission rate according to increase of channel transmission capacity can increase theoretically as much as a value obtained by multiplying a maximum transmission rate Ro, which corresponds to a case where one antenna is used, by an increase rate Ri of the following Equation 1.Ri=min(NT,NR)\u2003\u2003[Equation 1]\nFor example, in a MIMO communication system that uses four transmitting antennas and four receiving antennas, a transmission rate that four times bigger than that of a single antenna system can be obtained. After theoretical capacity increase of the MIMO system has been proved in the middle of 1990, various technologies have been actively studied to substantially improve a data transmission rate. Some of the technologies have been already reflected in the standard of various wireless communications such as third generation mobile communication and next generation wireless LAN.\nUpon reviewing the recent trend of studies related to the MIMO system, active studies are ongoing in view of various aspects such as the study of information theoretical aspect related to MIMO communication capacity calculation under various channel environments and multiple access environments, the study of radio channel measurement and model of a MIMO system, and the study of time space signal processing technology for improvement of transmission reliability and transmission rate.\nIn order to describe a communication method in a MIMO system in more detail, mathematical modeling of the communication method can be expressed as follows. As illustrated in FIG. 1, it is assumed that N7 transmitting antennas and NR receiving antennas exist. First of all, a transmitting signal will be described. If there exist N7 transmitting antennas, since the number of maximum transmission information is N7, the transmission information can be expressed by a vector shown in Equation 2 as follows.s==\u2514s1,s2, . . . ,sNT\u2518T\u2003\u2003[Equation 2]\nMeanwhile, different kinds of transmission power can be applied to each of the transmission information s1, s2, . . . , sNT. At this time, supposing that each transmission power is P1, P2, . . . , PNT, transmission information of which transmission power is controlled can be expressed by a vector shown in Equation 3 as follows.\u015d=[\u015d1,\u015d2, . . . ,\u015dNT]T=[P1s1,P2s2, . . . ,PNTsNT]T \nAlso, \u015d can be expressed by Equation 4 below using a diagonal matrix P.\n                    s        =                                            [                                                                                          P                      1                                                                                                                                                                                                                                                                                0                                                                                                                                                                                                                P                      2                                                                                                                                                                                                                                                                                                                                                                                                                                                                                                                    \u22f1                                                                                                                                                                                          0                                                                                                                                                                                                                                                                                  P                                              N                        T                                                                                                        ]                        \u2061                          [                                                                                          s                      1                                                                                                                                  s                      2                                                                                                            \u22ee                                                                                                              s                                              N                        T                                                                                                        ]                                =          Ps                                    [                  Equation          \u2062                                          \u2062          4                ]            \nMeanwhile, it is considered that a weight matrix W is applied to the information vector \u015d of which transmission power is controlled, so as to obtain N7 transmitting signals x1, x2, . . . , xNT. In this case, the weight matrix serves to properly distribute the transmission information to each antenna. Such transmitting signals x1, x2, . . . , xNT can be expressed by Equation 5 below using a vector X. In this case, Wij means a weight value between the i th transmitting antenna and the j th information. W may be referred to as a weight matrix or precoding matrix.\n                                                        x              =                            \u2062                              [                                                                                                    x                        1                                                                                                                                                x                        2                                                                                                                        \u22ee                                                                                                                          x                        i                                                                                                                        \u22ee                                                                                                                          x                                                  N                          T                                                                                                                    ]                                                                                        =                            \u2062                                                [                                                                                                              w                          11                                                                                                                      w                          12                                                                                            \u2026                                                                                              w                                                      1                            \u2062                                                          N                              T                                                                                                                                                                                                                    w                          21                                                                                                                      w                          22                                                                                            \u2026                                                                                              w                                                      2                            \u2062                                                          N                              T                                                                                                                                                                                          \u22ee                                                                                                                                                                                          \u22f1                                                                                                                                                                                                                                                            w                                                      i                            \u2062                                                                                                                  \u2062                            1                                                                                                                                                w                                                      i                            \u2062                                                                                                                  \u2062                            2                                                                                                                      \u2026                                                                                              w                                                      iN                            T                                                                                                                                                              \u22ee                                                                                                                                                                                          \u22f1                                                                                                                                                                                                                                                            w                                                                                    N                              T                                                        \u2062                            1                                                                                                                                                w                                                                                    N                              T                                                        \u2062                            2                                                                                                                      \u2026                                                                                              W                                                                                    N                              T                                                        \u2062                                                          N                              T                                                                                                                                                            ]                                \u2061                                  [                                                                                                                                          s                            ^                                                    1                                                                                                                                                                                          s                            ^                                                    2                                                                                                                                    \u22ee                                                                                                                                                                  s                            ^                                                    j                                                                                                                                    \u22ee                                                                                                                                                                  s                            ^                                                                                N                            T                                                                                                                                ]                                                                                                        =                            \u2062                              W                \u2062                                  s                  ^                                                                                                        =                            \u2062              WPs                                                          [                  Equation          \u2062                                          \u2062          5                ]            \nIf there exist NR receiving antennas, receiving signals y1, y2, . . . , yNR of the respective antennas can be expressed by a vector of Equation 6 below.y=[y1,y2, . . . ,yNR]\u2003\u2003[Equation 6]\nMeanwhile, in case of channel modeling in the MIMO communication system, channels can be classified depending on indexes of transmitting and receiving antennas. In this case, a channel that passes from the jth transmitting antenna to the ith receiving antenna will be expressed as hij. It is noted that index of the receiving antenna is prior to index of the transmitting antenna in index of hij.\nSeveral channels can be grouped into one and then can be expressed by a vector type or a matrix type. The channels of a vector type will be described below. FIG. 2 is a diagram illustrating a channel from N7 transmitting antennas to the ith receiving antenna.\nAs illustrated in FIG. 2, the channel from N7 transmitting antennas to the ith receiving antenna can be expressed as follows.hiT=\u2514hi1,hi2, . . . ,hiNT\u2518\u2003\u2003[Equation 7]\nAlso, all channels from N7 transmitting antennas to NR receiving antennas can be expressed by Equation 8 below through the matrix of Equation 7.\n                    H        =                              [                                                                                h                    1                    T                                                                                                                    h                    2                    T                                                                                                \u22ee                                                                                                  h                    i                    T                                                                                                \u22ee                                                                                                  h                                          N                      R                                        T                                                                        ]                    =                      [                                                                                h                    11                                                                                        h                    12                                                                    \u2026                                                                      h                                          1                      \u2062                                              N                        T                                                                                                                                                              h                    21                                                                                        h                    22                                                                    \u2026                                                                      h                                          2                      \u2062                                              N                        T                                                                                                                                          \u22ee                                                                                                                                                              \u22f1                    \u2062                                                                                                                                                                                                                                                                          h                                          i                      \u2062                                                                                          \u2062                      1                                                                                                            h                                          i                      \u2062                                                                                          \u2062                      2                                                                                        \u2026                                                                      h                                          iN                      T                                                                                                                    \u22ee                                                                                                                                          \u22f1                                                                                                                                                                                          h                                                                  N                        R                                            \u2062                      1                                                                                                            h                                                                  N                        R                                            \u2062                      2                                                                                        \u2026                                                                      h                                                                  N                        R                                            \u2062                                              N                        T                                                                                                                  ]                                              [                  Equation          \u2062                                          \u2062          8                ]            \nSince additive white Gaussian noise (AWGN) is actually added to the channels after the above channel matrix H, AWGN n1, n2, . . . , nNR added to each of the NR receiving antennas can be expressed by a vector of Equation 9 below.n=[n1,n2, . . . ,nNR]\u2003\u2003[Equation 9]\nThe receiving signals obtained using the above equations can be expressed by Equation 10 below.\n                                                        y              =                            \u2062                              [                                                                                                    y                        1                                                                                                                                                y                        2                                                                                                                        \u22ee                                                                                                                          y                        i                                                                                                                        \u22ee                                                                                                                          y                                                  N                          R                                                                                                                    ]                                                                                        =                            \u2062                                                                    [                                                                                                                        h                            11                                                                                                                                h                            12                                                                                                    \u2026                                                                                                      h                                                          1                              \u2062                                                              N                                T                                                                                                                                                                                                                                      h                            21                                                                                                                                h                            22                                                                                                    \u2026                                                                                                      h                                                          2                              \u2062                                                              N                                T                                                                                                                                                                                                          \u22ee                                                                                                                                                                                                                                      \u22f1                            \u2062                                                                                                                                                                                                                                                                                                                                                                                                  h                                                          i                              \u2062                                                                                                                          \u2062                              1                                                                                                                                                            h                                                          i                              \u2062                                                                                                                          \u2062                              2                                                                                                                                \u2026                                                                                                      h                                                          iN                              T                                                                                                                                                                            \u22ee                                                                                                                                                                                                          \u22f1                                                                                                                                                                                                                                                                                  h                                                                                          N                                R                                                            \u2062                              1                                                                                                                                                            h                                                                                          N                                R                                                            \u2062                              2                                                                                                                                \u2026                                                                                                      h                                                                                          N                                R                                                            \u2062                                                              N                                T                                                                                                                                                                          ]                                    \u2061                                      [                                                                                                                        x                            1                                                                                                                                                                            x                            2                                                                                                                                                \u22ee                                                                                                                                                  x                            j                                                                                                                                                \u22ee                                                                                                                                                  x                                                          N                              T                                                                                                                                            ]                                                  +                                  [                                                                                                              n                          1                                                                                                                                                              n                          2                                                                                                                                    \u22ee                                                                                                                                      n                          i                                                                                                                                    \u22ee                                                                                                                                      n                                                      N                            R                                                                                                                                ]                                                                                                        =                            \u2062                              Hx                +                n                                                                        [                  Equation          \u2062                                          \u2062          10                ]            \n(3) Definition of Reference Signal\nWhen a mobile communication system transmits packets, since the packets are transmitted through a radio channel, signal distortion may occur during transmission of the packets. In order to normally receive the distorted signal, a receiver should find out channel information and correct distortion of the signal as much as the channel information. In order to find out the channel information, when a signal known by both the transmitter and the receiver is received, the receiver finds out the channel information using a distortion status of the received signal. The signal known by both the transmitter and the receiver will be referred to as a pilot signal a reference signal (RS). The RS includes four types as illustrated in Table 1 depending on its purpose of use.\nTABLE 1FeaturesCQI-Common RS for channel measurement. Since UE(UseCRSEquipment) determines CQI(Channel Quality Indicator),PMI(Precoding Matrix Index), RI(Rank Indicator), etcbased on the CQI-CRS, the CQI-CRS is preferablydistributed over all bands.DM-Common RS for demodulation. But DM-CRS may be used forCRSchannel measurement. Since several UEs use the DM-CRSfor channel measurement, a precoding scheme of a specificUE cannot be applied to the DM-CRS. Accordingly, ifPDSCH(Physical Downlink Shared CHannel) is precoded,signaling is required through PDCCH(Physical DownlinkControl CHannel).NDM-Dedicated RS for non-precoded demodulation.DRSPDM-Dedicated RS for precoded demodulation. Since precodingDRSis equally applied to both PDM-DRS and PDSCH, signalingis separately required for PDCCH.\nRecently, when most of mobile communication systems transmit and receive data packets using multiple antennas, a separate reference signal for each of transmitting antennas exists, and should know a channel status between each transmitting antenna and each receiving antenna to normally receive a signal.\nMeanwhile, FIG. 3A is a diagram illustrating a structure of RS of normal cyclic prefix (CP) in a 3GPP LTE, and FIG. 3B is a diagram illustrating a structure of RS of extended cyclic prefix (CP) in a 3GPP LTE. In the current 3GPP LTE, a structure of a downlink RS is as illustrated in FIG. 3A and FIG. 3B. In one resource block, a horizontal axis means a time axis while a vertical axis means a frequency axis. One subframe includes two slots. FIG. 3A corresponds to a normal CP, wherein each slot includes seven symbols. FIG. 3B corresponds to extended CP, wherein each slot includes six symbols. The extended CP is generally used under the environment that delay is long. The number of base station transmitting antenna ports which are supported is four, and resources used to transmit RS for each of the antenna ports 0\u02dc3 are represented by 0, 1, 2 and 3.\n(4) Time/Frequency Spacing of RS\nIn the 3GPP LTE, when considering 20 MHz system, one symbol includes 2048 samples. In this case, symbol duration is 66.67 \u03bcs. 1 \u03bcs corresponds to 10\u22126 seconds. Meanwhile, in case of normal CP, 160 samples are used in the first one of seven symbols as CP, and 144 samples are used in the other six symbols as CP. At this time, each of the symbols can cover channel delay spread equivalent to 5.2 \u03bcs and 4.69 \u03bcs. In case of the extended CP, since six symbols include 512 samples, they can cover channel delay spread equivalent to 16.67 \u03bcs. As seen from FIG. 3A and FIG. 3B, a transmission interval between RS for same antenna ports on the frequency axis of one symbol is 6. However, since staggering is given within one slot, a subcarrier interval between RS on the frequency axis is finally 3. In the present invention, the interval will mean a subcarrier interval on the frequency axis unless specified otherwise. In this case, it is assumed that the subcarrier interval is 15 kHz.\nMeanwhile, the transmission interval of the RS and its channel estimation capability will be considered as follows using equations.\nFor example, if inverse discrete fourier transform is applied to sequence x[i]={x0, x1, x2, . . . , xN-1}, the following Equation 11 can be obtained.\n                                          X            \u2061                          [              k              ]                                =                                    \u2211                              i                =                0                                            N                -                1                                      \u2062                                          x                \u2061                                  [                  i                  ]                                            \u00b7                              \u2147                                  j                  \u2062                                                                          \u2062                  2                  \u2062                  \u03c0                  \u2062                                                                          \u2062                  k                  \u2062                                      \u2148                    N                                                                                      ,                  k          =          0                ,        1        ,        \u2026        \u2062                                  ,                  N          -          1                                    [                  Equation          \u2062                                          \u2062          11                ]            \nSince RS is sent between data at constant intervals, if 0 is filled in N number of data between the respective RSs, sequence\n      y    \u2061          [      i      ]        =      {                  x        0            ,                                    0            ,            0            ,            \u2026            \u2062                                                  ,            0                    \ufe38                          z          \u2062                                          \u2062          zeros                    ,              x        1            ,                                    0            ,            0            ,            \u2026            \u2062                                                  ,            0                    \ufe38                          z          \u2062                                          \u2062          zeros                    ,              x        2            ,      \u2026      \u2062                          ,              x                  N          -          1                    ,                                    0            ,            0            ,            \u2026            \u2062                                                  ,            0                    \ufe38                          z          \u2062                                          \u2062          zeros                      }  is obtained. If IDFT is applied to y[i], the following Equation 12 can be obtained.\n                                          Y            \u2061                          [              k              ]                                =                                    \u2211                              m                =                0                                                              N                  \u2061                                      (                                          Z                      +                      1                                        )                                                  -                1                                      \u2062                                          y                \u2061                                  [                  m                  ]                                            \u00b7                              \u2147                                  j                  \u2062                                                                          \u2062                  2                  \u2062                  \u03c0                  \u2062                                                                          \u2062                  k                  \u2062                                      m                                          N                      \u2061                                              (                                                  Z                          +                          1                                                )                                                                                                                                ,                                  \u2062                  k          =          0                ,        1        ,        \u2026        \u2062                                  ,                              N            \u2061                          (                              Z                +                1                            )                                -          1                                    [                  Equation          \u2062                                          \u2062          12                ]            \nIn the Equation 12, since y[m] and x[i] are\n  {                                                                        y                \u2061                                  [                  m                  ]                                            =                              x                \u2061                                  [                  i                  ]                                                      ,                                                              m              =                                                (                                      Z                    +                    1                                    )                                \u00b7                i                                      ,                                                              i              =              0                        ,            1            ,            \u2026            \u2062                                                  ,                          N              -              1                                                                                                      y                \u2061                                  [                  m                  ]                                            =              0                        ,                                                otherwise            ,                                                                                \u2003  the relationship between Y[k] X[k] can be observed using the Equations 11 and 12 in such a manner that signals are repeatedly generated as much as the number of 0 as illustrated in Equation 13.\n                              Y          \u2061                      [            k            ]                          =                  {                                                                                          X                    \u2061                                          [                      0                      ]                                                        ,                                      X                    \u2061                                          [                      1                      ]                                                        ,                  \u2026                  \u2062                                                                          ,                                      X                    \u2061                                          [                                              N                        -                        1                                            ]                                                                      \ufe38                                            N                \u2062                                                                  \u2062                samples                                      ,            \u2026            \u2062                                                  ,                                                                                X                    \u2061                                          [                      0                      ]                                                        ,                                      X                    \u2061                                          [                      1                      ]                                                        ,                  \u2026                  \u2062                                                                          ,                                      X                    \u2061                                          [                                              N                        -                        1                                            ]                                                                      \ufe38                                            N                \u2062                                                                  \u2062                samples                                              }                                    [                  Equation          \u2062                                          \u2062          13                ]            \nFor example, if the interval of RS is 6, six signals are repeatedly generated in accordance with the Equation 13 as illustrated in FIG. 4. In this case, if channel delay spread does not exceed 11.11 \u03bcs, signal recovery can be performed theoretically. Namely, in case of the normal CP, since the number of CP samples is 160, 144, signal recovery can be performed even in the case that the transmission interval of RS is 6 when channel delay spread of maximum 5.2 \u03bcs and 4.69 \u03bcs occurs.\nMeanwhile, in case of the extended CP, since the number of CP samples is 512, it is supposed that channel delay spread is maximum 16.67 \u03bcs. Accordingly, a problem occurs in that interference exists between symbols as illustrated in FIG. 5.\nIn order to solve the aforementioned problem, in the LTE, staggering is applied to RS of each antenna port, so that the transmission interval of the RS becomes 3. As illustrated in FIG. 6, if the transmission interval of the RS becomes 3, channel delay spread that can be covered reaches 22.22 \u03bcs, whereby signal recovery can be performed regardless of the length of the CP.\n(5) Cyclic Delay in Time Domain\nIf sequence x[i]={x0, x1, x2, . . . , xN-1} is cyclic shifted on the time axis, the relationship with the frequency axis will be observed as illustrated in Equation 4 below.\n                                                                                                                                    X                      \u2032                                        \u2061                                          [                                              k                        \u2032                                            ]                                                        =                                    \u2062                                                            \u2211                                              i                        =                        0                                                                    N                        -                        1                                                              \u2062                                                                                            x                          \u2061                                                      [                            i                            ]                                                                          \u00b7                                                  \u2147                                                      j                            \u2062                                                                                                                  \u2062                            2                            \u2062                            \u03c0                            \u2062                                                                                          k                                \u2032                                                            N                                                        \u2062                            \u2148                                                                                              \u2062                                              \u2147                                                  j                          \u2062                                                                                                          \u2062                          2                          \u2062                          \u03c0                          \u2062                                                      u                            M                                                    \u2062                          \u2148                                                                                                                    ,                                                                                                        \u2062                                                      N                    =                    MZ                                    ,                                                            k                      \u2032                                        =                    0                                    ,                  1                  ,                  \u2026                  \u2062                                                                          ,                                      N                    -                    1                                                                                                                          =                                \u2062                                                      \u2211                                          i                      =                      0                                                              N                      -                      1                                                        \u2062                                                            x                      \u2061                                              [                        i                        ]                                                              \u00b7                                          \u2147                                                                        j2\u03c0                          \u2061                                                      (                                                                                                                            k                                  \u2032                                                                +                                uZ                                                            N                                                        )                                                                          \u2062                        \u2148                                                                                                                                \u2062                                  \u2062                                            X              \u2032                        \u2061                          [                                                (                                      k                    -                    uZ                                    )                                \u2062                mod                \u2062                                                                  \u2062                N                            ]                                =                                    \u2211                              i                =                0                                            N                -                1                                      \u2062                                          x                \u2061                                  [                  i                  ]                                            \u00b7                              \u2147                                  j2\u03c0                  \u2062                                      k                    N                                    \u2062                  \u2148                                                                    \u2062                                  \u2062                                            X              \u2032                        \u2061                          [                                                (                                      k                    -                    uZ                                    )                                \u2062                mod                \u2062                                                                  \u2062                N                            ]                                =                      X            \u2061                          [              k              ]                                                          [                  Equation          \u2062                                          \u2062          14                ]            \nIn other words, if\n  \u2147      j2\u03c0    \u2062          u      M        \u2062    \u2148  is multiplied on the frequency axis, cyclic delay can occur on the time axis.\nIf the number of transmitting antennas that can be supported for a system that includes N number of antennas is additionally provided to support M(M>N) number of transmitting antennas, there exist an LTE UE (user equipment) that can recognize the existing N number of transmitting antennas only and an LTE-A UE that can recognize M number of transmitting antennas only. In this case, in addition to RS for supporting the existing N number of transmitting antennas, transmission of additional M-N number of RSs is required. At this time, under the environment that the LTE-A UE that recognizes M number of antennas is added to the LTE UE that recognizes the existing N number of antennas without additional signaling, efficient transmission of data and RS is required."}
-{"text": "In a digital video broadcast system, a headend component uses Quadrature Amplitude Modulator (QAM) devices or Internet Protocol (IP) services to deliver audio/video (AV) content to a set-top box as a Moving Pictures Expert Group (MPEG) stream. The MPEG stream may include \u201clive\u201d advertisements or targeted advertisements. The targeted advertisements are placed in a parallel QAM or an IP service that is synchronized with the \u201clive\u201d advertisements. Alternatively, the targeted advertisements may be stored on a data storage device that is either integrated with the set-top box, or accessible via a network communication connection. Software on the set-top box that includes a digital video recorder (DVR) determines whether to perform a switch to the targeted advertisement based on various traits, such as demographics, purchase history, observed behavior, or the like.\nMultiple System Operator (MSO) customers are developing and deploying targeted advertising systems. These systems enable a set-top box to switch from a network service to an addressable advertising service, and back to the network service during an advertising break. The term \u201cswitch\u201d includes the definition of \u201cswitch\u201d in the context of the American National Standard (ANSI) Society of Cable Telecommunications Engineers (SCTE) Standard ANSI/SCTE 138 (i.e., switch describes tunes (L0 switch) and splices (L1 switch)), the definition of \u201ctune\u201d in the context of a QAM device, and the definition of \u201cswitch\u201d or \u201ctransition\u201d in the context of an IP or DVR device. In the first phase of deployment, the delivery of the network service and addressable advertising services are on separate transport streams, in which the addressable advertising content is synchronized with an advertisement on the network service. Since multiple advertisement streams are simultaneously delivered to the set-top box, the set-top box selects one of the advertisements, based on received information and locally stored selection criteria, and switches at the appropriate time.\nThe addressable advertisements contain filler, as defined in ANSI/SCTE 138. The filler follows signaled switch points which mark an opportunity to switch from the network service to the addressable advertising service, and an opportunity to switch from the addressable advertising service back to the network service. The filler is properly formatted frames that represent black pictures and silent audio, and that optimize the service acquisition procedure by allowing the set-top box the time necessary to switch to and begin decoding the new service. The duration of the filler is selected to allow for the time needed for the slowest set-top box in the network to switch to and begin decoding the new service. Since the duration of the filler is lost on the switch to the addressable advertising service and on the switch back to the network service, the duration of the filler takes away from the duration of the addressable advertising content. There is a need to minimize the duration of the filler and maximize the duration for the addressable advertising content. The presently disclosed invention satisfies this demand."}
-{"text": "A gas turbine is conventionally comprised of a compressor, a combustor, and a turbine. The turbine is coupled to the compressor in order to drive the compressor. The combustion chamber receives fuels such as a combustion gas, and a certain amount of nitrogen, to lower the flame temperature in the combustion chamber of the combustor, which makes it possible to minimize the discharge of nitrogen oxides to atmosphere. The combustion gas may be obtained by gasification, that is, oxidation of carbon products such as coal. This partial oxidation is carried in an independent unit referred to as a gasifier. Conventionally, the gas turbine is combined with an air separation unit. The air separation unit enables at least one gas stream, mostly consisting of one of the gases of air, especially oxygen or nitrogen, to be supplied from input air. To combine the air separation unit with the gas turbine, the oxygen and nitrogen produced in the air separation unit are admitted respectively into the gasifier and the combustion chamber of the combustor."}
-{"text": "1. Field of the Invention\nThe present invention relates to a method of operating a split gate type of non-volatile memory cell and a semiconductor memory device having the cells, and more particularly, to a method of operating a split gate type of non-volatile memory cell overcoming the problems of the program disturbance and endurance characteristics and a method of operating a semiconductor memory device including the cells.\n2. Discussion of Related Art\nThe split gate type of non-volatile memory cell is known in U.S. Pat. No. 5,045,488 entitled xe2x80x9cmethod of manufacturing a single transistor non-volatile semiconductor devicexe2x80x9d and U.S. Pat. No. 5,029,130 entitled xe2x80x9csingle transistor non-volatile alterable semiconductor memory devicexe2x80x9d. Methods of programming memory cells having floating gate electrodes are disclosed in U.S. Pat. No. 5,659,504 to Bude et al.\nFIG. 1 illustrates the structure of a conventional split gate type of non-volatile memory cell in U.S. Pat. Nos. 5,045,488 and 5,029,130.\nA source 12 and a drain 14 are formed on a semiconductor substrate 10. Between the source 12 and drain 14 is formed a channel 16. An insulating layer 18 is formed on the source 12, channel 16 and drain 14. A floating gate 20 is formed on a predetermined portion of the insulating layer 18 on the channel 16 and drain 12. Other insulating layer 22 is formed on the floating gate 20. Another insulating layer 24 is formed to be insulated from a control gate 26. The control gate 26 is formed on a predetermined portion of the insulating layer 24 and the insulating layer 18 on the source 12 and channel 16.\nThe operation of the split gate type of non-volatile memory cell of FIG. 1 is described with reference to FIGS. 2 to 4.\nFIG. 2 illustrates a method of erasing the split gate type of non-volatile memory cell of FIG. 1. In the drawing, the source 12 and drain 14 receive the same voltage 0V, and the control gate 26 receives the voltage Vpp higher than that applied to the source 12 and drain 14. Here, an intensive coupling from the floating gate 20 to the substrate 10 and drain 14 decreases a voltage of the floating gate 20. This kind of voltage decrease allows electrons to flow from the floating gate 20 to the control gate 26 by an F-N (Fowler-Nordheim) tunneling mechanism. Accordingly, an erasing function is ascribed to the electron of the floating gate 20 moving to the control gate 26. Through the erasing operation, the floating gate 20 is charged with (+). That is, the erasing operation is performed by the voltage difference between the floating gate 20 and the control gate 26.\nFIG. 3 illustrates a method of programming the split gate type of non-volatile memory cell of FIG. 1. A threshold voltage Nth is applied to the control gate 26. A high voltage Vpp is applied to the drain 14. xe2x80x9c0xe2x80x9d voltage is applied through the source 12 to the substrate 10, and thus the programming is performed.\nIf the high voltage Vpp is applied to the drain, the potential of the floating gate 20 is raised and the channel under the floating gate 20 is turned on. The threshold voltage Vth is applied to the control gate 26 and accordingly the channel under the control gate 26 is lightly turned on. Accordingly, electrons flow from the source 12 to drain 14. These electrons are charged in the floating gate 20 via the insulating layer 18 because of the static electricity of the floating gate 20, thus performing the programming operation. Hence, the floating gate 20 is (xe2x88x92)-charged and programmed to xe2x80x9c0xe2x80x9d.\nIn other words, the programming operation is performed in such a manner that a high voltage Vpp is applied to the drain of the memory cell to thereby bring the floating gate 20 to a predetermined voltage, and a predetermined voltage (a threshold voltage Vth of a transistor made of the control gate and the channel) is applied to the control gate 26 so that hot channel electrons generated when the current flows between the source 12 and drain 14 are injected into the floating gate.\nFIG. 4 illustrates a method of reading the conventional split gate type of non-volatile memory cell of FIG. 1. A reference voltage Vref is applied to the control gate 26; 0V to the drain 14, 2V to the source 12, and 0V to the substrate 10, thus performing the reading operation.\nIf the floating gate 20 is charged with (+), the channel 16 right under the floating gate 20 is turned on. If the voltage of the control gate 26 is raised to the reference voltage Vref for turning on the channel under the control gate 26, electrons may flow from the drain 14 to the source 12, and thereby reading the data of xe2x80x9c1xe2x80x9d.\nOn the contrary, if the floating gate 20 is charged with (xe2x88x92), the channel right under the floating gate 20 is slightly turned on or off. A voltage level of the control gate 26 and the source 12 is raised to that of the read voltage so as to turn on the channel under the control gate 26. Hence, the current cannot flow through the channel, thereby reading the data of xe2x80x9c0xe2x80x9d.\nIn other words, if the floating gate 20 is charged with (+), the current is generated through the channel 16 to thereby read the data of xe2x80x9c1xe2x80x9d, and on the contrary, if the floating gate 20 is charged with (xe2x88x92), the current does not flow through the channel 16, thereby reading the data of xe2x80x9c0xe2x80x9d.\nTherefore, the data is read by checking if the current flows through the memory cell or not by applying a predetermined voltage to the source 14 and the control gate 26. For this reason, if to perform the reading operation, the channels should be formed through the control gate and floating gate so that the current flows through the cell.\nBut, in the semiconductor memory device including the thus-structured split gate type of non-volatile memory cells, a threshold voltage Vth should be applied to the word line of a selected cell; 0V to the word line of a non-selected cell; 0V to the bit line of a selected cell; power voltage Vcc to the bit line of a non-selected cell; high voltage Vpp to the drain of a selected cell; and 0V to the drain of a non-selected cell should be respectively applied in order to generate current between the source and drain, thereby performing the programming operation. Within the programming condition, if the non-selected cell including the drain commonly connected to the drain of a selected cell is erased, the high voltage Vpp is applied to the drain of a non-selected cell, the floating gate being charged with (+), the source receiving 0V, and the substrate receiving 0V. Here, the control gate and source receive 0V, but the channel is formed by the punch through phenomenon and the current flows. The electrons conducted in the channel are injected into the floating gate, and thus programmed. Consequently, an ON-cell becomes an OFF-cell, thus causing the program interference problem.\nIn addition, the semiconductor memory device including the conventional split gate-typed memory cells repeatedly performs the programming and erasing operations. In the erasing operation, electrons of the floating gate should completely go out towards the control gate. But, they are trapped within the tunneling insulating layers. Consequently, the threshold voltage increases in accordance with the increasing number of operations, thus causing the problem of endurance characteristics.\nIt is therefore an object of the present invention to provide improved methods of operating non-volatile memory devices.\nIt is another object of the present invention to provide methods of operating non-volatile memory devices that enhance the endurance characteristics of memory cells therein.\nIt is still another object of the present invention to provide methods of operating non-volatile memory devices that increase the reliability and speed of programming and erasing operations.\nIt is still a further object of the present invention to provide methods of operating non-volatile memory devices that reduce the likelihood of programming interference between selected and non-selected memory cells.\nTheses and other objects, advantages and features of the present invention are provided by preferred methods of operating non-volatile memory cells (e.g., EEPROM cells) that include the use of negative substrate biases during programming and erasing operations. In particular, a preferred method of operating a non-volatile memory cell includes the step of erasing the memory cell by withdrawing negative charge from a floating gate therein using a positive control electrode bias and a negative substrate bias. The use of a negative substrate bias increases the potential difference between the control electrode and the floating gate and this increase results in faster and more reliable erasing. The preferred method also includes the step of programming the memory cell by accumulating negative charge on the floating gate using a positive control electrode bias, a negative substrate bias and a positive drain bias. Here, the negative substrate bias is used advantageously to reduce the likelihood that non-selected memory cells will become inadvertently programmed during operations to program selected memory cells."}
-{"text": "A plurality of radio communication systems such as a cell-phone system and a radio MAN (Metropolitan Area Network) are currently used. For attaining a further speeding up and large capacity of radio communication, lively discussion is continuously performed about a next generation radio communication technology.\nFor example, in a 3GPP (3rd Generation Partnership Project) being a standardization organization, there is proposed a communication standard referred to as an LTE (Long Term Evolution) enabling communication using a frequency band of 20 MHz at a maximum. Further, as a next generation communication standard of LTE, there is proposed a communication standard referred to as an LTE-A (LTE-Advanced) enabling communication using five frequency bands (namely, a frequency band of 100 MHz) of 20 MHz at a maximum (see, for example, Non-Patent Literatures 1 and 2). In the LTE-A, the number of frequency bands to be used is proposed to be dynamically changed according to traffic (see, for example, Non-Patent Literature 3).\nFurther, in a radio communication system, from one radio communication device (e.g., a mobile station) to another radio communication device (e.g., a base station) which performs allocation control of radio resources, a random access may be performed. The random access from the mobile station to the base station is performed, for example, at the time when (1) the mobile station first accesses the base station, (2) an allocation of radio resources used for data transmission is requested to the base station, and (3) synchronization is established during reception of data from the base station, and (4) synchronization is achieved with a mobile target base station during a handover.\nThe random access includes a contention based random access and a non-contention based random access (see, for example, 10. 1. 5 section of Non-Patent Literature 4, and 5. 1 section of Non-Patent Literature 5). In the case of the random access from the mobile station to the base station, in the contention based random access, the mobile station arbitrarily selects a signal sequence from among a plurality of signal sequences and transmits it to the base station as a random access preamble. In the non-contention based random access, the base station notifies the mobile station of information in which a signal sequence is specified and the mobile station transmits a signal sequence according to the notification from the base station as the random access preamble.\nNPTL1: 3GPP (3rd Generation Partnership Project), \u201cRequirements for further advancements for Evolved Universal Terrestrial Radio Access (E-UTRA) (LTE-Advanced)\u201d, 3GPP TR 36.913 V8.0.1, 2009-03.\nNPTL2: 3GPP (3rd Generation Partnership Project), \u201cFeasibility study for Further Advancements for E-UTRA (LTE-Advanced)\u201d, 3GPP TR 36.912 V9.0.0, 2009-09.\nNPTL3: 3GPP (3rd Generation Partnership Project), \u201cThe need for additional activation procedure in carrier aggregation\u201d, 3GPP TSG-RAN WG2 #67bis R2-095874, 2009-10.\nNPTL4: 3GPP (3rd Generation Partnership Project), \u201cEvolved Universal Terrestrial Radio Access (E-UTRA) and Evolved Universal Terrestrial Radio Access Network (E-UTRAN); Overall description\u201d, 3GPP TS 36.300 V9.0.0, 2009-06.\nNPTL5: 3GPP (3rd Generation Partnership Project), \u201cEvolved Universal Terrestrial Radio Access (E-UTRA) Medium Access Control (MAC) protocol specification\u201d, 3GPP TS 36.321 V9.1.0, 2009-12.\nIncidentally, in a radio communication system capable of performing communication by using a plurality of frequency bands, the number of frequency bands to be used according to traffic as described above is considered to be changed. However, in a method as described in the Non-Patent Literature 3, after communication is started between radio communication devices (after completing a random access procedure), a procedure is freshly performed so as to use other frequency bands except the frequency band in which communication is started. In this method, in the case where it is proved that the other frequency bands are desired to be used before starting communication (for example, in the case where a transmission data amount is proved to be large), the procedure becomes inefficient."}
-{"text": "This invention relates to electrical connectors, and, more particularly, to separable electrical connectors suited for use under high-voltage conditions. Still more particularly, this invention relates to gas actuated high-voltage bushings having a contact mounted within a bore for reciprocal movement within a bushing housing.\nHigh-voltage separable connectors innerconnect sources of energy such as transformers to distribution networks or the like. The situations typically encountered in the connection and disconnection of electrical connectors and power distributions include \"loadmake\", \"loadbreak\", and \"fault closure\". Loadmake includes the joinder of male and female contact elements, one energized and the other engaged with a normal load. An arc of moderate intensity is struck between the contact elements as they approach one another and until joinder. Loadbreak includes the separation of such mated male and female contact elements, while they supply power to a normal load. Moderate intensity arcing again occurs between the contact elements from the point of separation thereof until they are somewhat removed from one another. Fault closure includes the joinder of male and female contact elements, one energized and the other engaged with a load having a fault, e.g., a short circuit condition. A substantial arcing occurs between the contact elements as they approach one another and until joinder, giving rise to the possibility of explosion and accompanying hazard to operating personnel.\nThe prior art teaches the use of materials which emit arc-quenching gas when subjected to arcing, thus adequately dissipating the moderate intensity of arcs which occur during loadmake and loadbreak. The problem situation is fault closure, in which considerably more arc-quenching gas and mechanical assistance are required to extinguish the arc. During fault closure, the gas generated pressures may be fifty times greater than such pressures during loadmake. With respect to fault closure, the prior art has relied upon the use of the arc-quenching gas to assist in accelerating the contact elements into engagement, thus minimizing arcing time.\nA typical prior art electrical connector includes a bushing well connected to the transformer, a bushing insert which contains a female contact assembly connected to the well, and an elbow connected to a distribution line and containing a male contact to join an insert female contact in the female contact assembly. Because closure of the male and female contacts can occur under activated conditions or under fault conditions, the female contact is arranged to move within the insert to hasten the closure of the male and female contacts and thus extinguish any arc created. However, it is necessary to maintain electrical continuity during the travel of the female contact assembly. The connection between such female contact assembly and the remainder of the bushing insert must be flexible so as not to impede its movement but sufficient to carry the high currents in the circuit. Typical prior art devices include a female contact which has a piston that is moveable between a first and second position. Gas pressure which is generated by arcing during fault closure accelerates the female contact toward the male contact, thus hastening contact engagement and decreasing the time duration of the are. Mechanisms for achieving these results have not always produced sufficient current paths causing the connectors to run hot, and interfering with proper operation of the distribution network and in the extreme, leading to the destruction of the bushing inserts."}
-{"text": "The present invention relates to an image display apparatus such as a head mounted display (HMD), and especially to an image display apparatus which forms images in mutually different viewing angle areas in order to display a combined image of the images joined to each other.\nObservation optical systems used in HMD are desired to be compact and lightweight. The observation optical systems are also desired to be capable of displaying an image with a wide viewing angle to increase realism in the displayed image.\nAmong observation optical systems for the HMD, a so-called ocular optical system type observation optical system is comparatively easy to be miniaturized and provided with a wide viewing angle. However, such an ocular optical system type observation optical system requires, for providing a wide viewing angle with one image-forming element, a large-size image-forming element having a large image-forming area where an original image is formed.\nIn contrast thereto, a so-called tiling technology is proposed which enables displaying of an image with a wide viewing angle by using plural small-size image-forming elements having a small image-forming area. The tiling technology joins plural enlarged images to each other which are formed by light fluxes from plural original images formed on the plural image-forming elements to enable the displaying of one enlarged combined image. Such observation optical systems (HMD) for displaying an image with a wide viewing angle by the tiling technology are disclosed in Japanese Patent Laid-Open No. 09-166759 and Japanese Patent Laid-Open No. 11-326820.\nJapanese Patent Laid-Open No. 09-166759 discloses an observation optical system in which two optical units are provided for one eye, each optical unit including an image-forming element and a three-surface prism having an entrance surface, an internal total reflection/exit surface and a reflection surface, and the two optical units are disposed at upper and lower positions so as not to be overlapped with each other in a right and left direction. In this observation optical system, two enlarged images are formed by the two optical units at right and left positions, and thereby one enlarged combined image with a horizontally wide viewing angle is displayed. When a cross section where a light flux (optical path) is folded by at least twice reflections by an optical element such as the three-surface prism is defined as a local meridional cross section, the \u201cright and left (horizontal) direction\u201d in Japanese Patent Laid-Open No. 09-166759 corresponds to a direction along a local sagittal cross section orthogonal to the local meridional cross section.\nJapanese Patent Laid-Open No. 09-166759 also discloses an example in which the two optical units disposed at the upper and lower positions are relatively moved so as to be overlapped with each other in a visual axis direction of an observer. A region where the two optical units are overlapped with each other allows observation of an outside optical image (so-called optical see-through observation). In this example, enlarged images formed by the two optical units are overlapped with each other at a central region in the right and left direction to form one enlarged combined image.\nJapanese Patent Laid-Open No. 11-326820 discloses an observation optical system including two division optical systems which form two enlarged images, the two division optical systems using one prism having a symmetric shape in a local meridional cross section. In this observation optical system, a light flux from an original image formed on one image-forming element enters the prism through an entrance surface thereof, is totally reflected by an internal total reflection/exit surface thereof, is then reflected by a concave mirror surface, and finally exits from the prism through the internal total reflection/exit surface to proceed to an exit pupil. A light flux from an original image formed on the other image-forming element proceeds to the same exit pupil in a same manner as that of the light flux from the one image-forming element. The two concave mirror surfaces that cause the light fluxes from the two image-forming elements to proceed to the same exit pupil are closely adjacent (joined) to each other.\nHowever, the latter observation optical system disclosed in Japanese Patent Laid-Open No. 09-166759 does not contribute to miniaturization of the observation optical system (in other words, of the HMD) in a direction of the local meridional cross section (vertical direction) in which the observation optical system is decentered, because its viewing angle is not divided in the direction of the local meridional cross section while the viewing angle is divided in a direction of the local sagittal cross section (horizontal direction) to obtain a wide viewing angle. Each of the observation optical systems disclosed in Japanese Patent Laid-Open No. 09-166759 aims to enable the see-through observation through the overlapped region of the two optical units. Therefore, in the local meridional cross section, the viewing angles of the two optical units are the same, and parts thereof are overlapped with each other. Therefore, image division cannot be performed in the local meridional cross section, which is not effective for reducing the thickness or vertical size of the observation optical system.\nMoreover, the former observation optical system disclosed in Japanese Patent Laid-Open No. 09-166759 and the observation optical system disclosed in Japanese Patent Laid-Open No. 11-326820 form a joined portion of the two images formed by the light fluxes from the two original images (corresponding to a joined portion of the two concave mirror surfaces) in the enlarged combined image. The joined portion of the two concave mirror surfaces generates light scattering and flare, which makes the joined portion of the images noticeable."}
-{"text": "1. Field of the Invention\nThe invention relates generally to the field of electronic circuits and more particularly to designs for SRAM memory cells that provide improved stability in comparison to conventional designs.\n2. Related Art\nComputer systems and other devices typically need to have means for storing information.\nThese means may include persistent storage devices for large amounts of data, as well as smaller memory systems for storing data that the computer or other device is currently using. The memory systems for storing currently used data include both read-only memory (ROM) and random access memory (RAM.)\nRAM is typically used as the working memory of a device. RAM is used by devices to store data that needs to be accessible by a processor, and also needs to be modifiable. That is, the data can be changed. By contrast, data stored in a ROM cannot be modified, but can only be read. There is a great demand for RAM in computers and other electronic devices because the more RAM a device has, the more data can be readily accessible to the device's processor. For example, in a computer, the availability of more RAM enables the computer to execute more (or larger) software applications without having to swap data between RAM and a persistent storage device, such as a hard disc drive.\nThere are various different types of RAM. For example, dynamic RAM, or DRAM, has often been used in computers. The \u201cdynamic\u201d aspect of DRAM refers to the fact that DRAM memory cells need to be periodically refreshed in order to maintain the data which is stored in the cells. If the DRAM cells are not refreshed, the data will be lost. Static RAM, or SRAM, is another type of memory that is often used in computers. That \u201cstatic\u201d aspect of SRAM refers to the fact that SRAM cells do not have to be refreshed in the same manner as DRAM cells.\nSRAM memory has a number of advantages over DRAM memory. As noted above, SRAM cells do not have to be refreshed in order to maintain the data that is stored in them. Additionally, SRAM is typically much faster than DRAM. For example, typical SRAM cells may have access times of about less than 1 nanosecond, while DRAM cells may have access times closer to 60 nanoseconds. Further, SRAM memories do not require pauses between accesses, so the cycle time to access SRAM cells is typically much shorter than the cycle time for accessing DRAM cells.\nAlthough it has a number of advantages, SRAM memory also has some disadvantages. For instance, SRAM is typically much more expensive to manufacture than DRAM. Because the cost of SRAM memory is much higher than DRAM, it is common for SRAM memory to be used as a memory cache, while DRAM is used for a processor's main memory.\nSRAM memory cells may also be unstable. That is, the data in the cells may actually be corrupted when the cells are read. This problem arises from the fact that SRAM cells are read by coupling the cells to pre-charged bit lines and allowing the cells to pull down these bit lines. In other words, a higher voltage on the bit line is coupled to a lower voltage in the SRAM cell, causing the bit line voltage to drop and the SRAM cell voltage to rise. The voltage drop on the bit line is detected and amplified to provide the data for the associated processor.\nThe voltage rise in the SRAM cell, however, may corrupt the data stored in the cell. (because the initial, low voltage corresponded to the \u201c0\u201d stored in the cell, and the higher voltage resulting from the access may cause the data to be ambiguous, or even to flip-flop, so that it is now a \u201c1.\u201d)\nIf the transistors that make up the components of the memory cell were exactly identical to each other, the probability that the data in the cell would be corrected in this manner would be relatively low. As a practical matter, however, the transistors are not identical, but instead have slight variations that cause them to have slight variations in their respective responses. For example, each transistor has an associated threshold voltage that affects the response of the transistor. Because of manufacturing differences between the transistors, there is a variation in these threshold voltages that, in turn, results in a variation in the transistor's responses. These variations are becoming increasingly significant because the variation in threshold voltages increases as the size of the transistors decreases. Thus, as the memory cells grow smaller, they are more susceptible to the instability problem.\nThe instability of SRAM cells is obviously problematic. It would clearly be desirable to provide a design for an SRAM cell that is more stable than conventional designs and is therefore less likely to be corrupted when the cell is read."}
-{"text": "1. Field of the Invention\nThe present invention relates to adaptive optics, and more specifically, it relates to a reference-free compensated imaging system for recovering a high resolution image of a target.\n2. Description of Related Art\nThe basic elements of a typical (prior art) adaptive optics system 100 are shown in FIG. 1. The goal of such systems is to provide real-time compensation for propagation errors, as encountered by optical beams, as they travel through dynamic distorting paths, including turbulent atmospheres, optical pointing errors, imperfect optical elements, multi-mode optical fibers, etc.\nBy compensating for optical wavefront distortions, one can enhance the performance of a variety of optical systems. Examples include optical communication systems, remote sensors, precision laser-beam delivery systems for industrial and medical purposes, and compensated imaging systems such as in medical applications (ophthalmological imaging and precision surgical procedures through the eye) and microscopy. In the latter example, this implies that one can view complex objects over a distorted path with the same image quality as if the path were distortion-free. In this case, the performance of the imaging system can approach that of its theoretical diffraction limit, within the so-called isoplanatic cone.\nIn what follows, we first discuss a generic adaptive optical system capable of correcting for path distortions encountered by a so-called reference beam. The reference beam is typically an image-free optical source, whose function is to sample the path distortions and, thus, provide this wavefront information as input to the adaptive optical system. This discussion is followed by a description of a specific adaptive optical configuration typical of prior-art, including an example of a wavefront-error sensing device. This, in turn, is followed by a discussion of an optical compensated imaging system typical of the art. An understanding of these prior-art systems will provide perspective with regard to the exemplary embodiments of this invention that follow.\nAs we discuss below, compensation of wavefront phase errors enables a system to provide diffraction-limited imaging and viewing of an extended object. In general, one first samples and compensates for propagation-path errors using, a diffraction-limited reference beam. Upon compensation of wavefront errors encountered by the reference beam, the optical system can approach its theoretical diffraction-limited imaging capability of image-bearing beams that lie within the so-called isoplanatic patch, which is well known in the art.\nIt is to be appreciated that the compensated optical imaging system can be implemented to service a variety of imaging based applications beyond atmospheric viewing systems. Hence, when the basic imaging system is referred to as a telescope, it is to be understood that the present teachings and embodiments can also be applied, without loss of generality, to compensated microscopy systems, speckle imaging, ophthalmological systems, communications systems, and the distortion path is referred to as a dynamic atmosphere, ultrasound imaging systems and so on. Similarly, when the distortion path that imposed the wavefront distortions to be compensated is referred to as a dynamic atmosphere, it is to be understood that the teachings can also be applied, without loss of generality, to correct for propagation-path distortions such as those experienced by imperfect optical elements, and static and/or dynamic distortions due to propagation through ocular systems, skin tissue, clouds, turbid liquids, and so on. The scene-based (Shack-Hartman) wave-front sensor could also be used in a post-processing scheme such as deconvolution or to augment speckle imaging.\nTurning now to FIG. 1, the goal of the prior art system is to enable one to view an optical source 110 with diffraction-limited capability. In this case, the optical source is chosen to be of spatial extent less than, or equal to, the diffraction limit of the optical system. Therefore, this source is equivalent to a point object with zero image-bearing information, analogous to a single pixel in an image. Light that emerges from this object, which is referred heretofore as a \u201creference beam,\u201d 120, propagates through space, and, in general, becomes aberrated, as depicted by wavefront 120, as a result of the intervening path distortions or spatial phase errors, labeled by \u03c6. In essence, the reference beam 120 samples the propagation path distortions between it and the optical compensation system, 100, including distortions imposed by optical elements within the compensation system itself.\nAt the receiver end of the link, a fraction of reference beam 120 is collected by telescope 130, which represents the input optical imaging elements of the adaptive optical receiver system 100. The collected light forms an image at the camera, or detector array, 190. In the absence of path distortions, the image at the camera plane would be in the form of an Airy disc, since the reference beam 120 is a sub-diffraction-limited point-source. However, owing to optical propagation phase distortions, \u03c6, encountered by the reference beam on its path toward the receiver 110, the wavefronts of this beam will be aberrated, resulting in a distorted image of an Airy disc pattern at camera 190. As is known in the art, the path distortions in this scenario can stem from atmospheric turbulence, pointing and tracking errors, imperfect optical elements, thermal and mechanical perturbations, among other effects. The goal, therefore, of the adaptive optical system 100 is to compensate for such path errors so that the image quality of the reference beam at detector 190 can approach the diffraction limit.\nReturning to FIG. 1, the reference beam exiting the telescope 130 will be aberrated by virtue of the deleterious path distortions, as represented by wavefront 140. In this example, the adaptive optical system consists of two optical correction elements. The first corrective element 150 is a so-called tip-tilt compensator, whose function is to compensate for overall beam pointing and tracking errors. The second corrective element 160 is a spatial phase modulator, whose function is to compensate for fine-scale optical phase errors, including focus errors and spatially complex wavefront errors. The latter can include static and/or dynamic errors resulting from atmospheric turbulence and surface and volume refractive-index irregularities of optical elements, as discussed above. Wavefront compensation element 160 can be in the form of arrays of continuous and/or discrete optical phase shifters, such as piezoelectric transducers, electro-optic elements, deformable membranes, MEMS mirrors, liquid crystal cells, photonic crystals, among other devices, as is known in the art.\nThe incident distorted beam 140, first encounters the tip-tilt optical component 150 followed by the spatial phase modulator 160. The beam subsequently strikes a beam splitter 165, with one output beam directed to an optical wavefront error sensor 170, and with the other output beam directed to the camera detector 190.\nThe telescope provides an image of the incident beam at the camera plane 190, and, furthermore, provides an image of the pupil plane at the surface of the wavefront corrective element 160. Hence, the wavefront at the incident aperture is replicated, and scaled, as needed, at the plane of 160. The number of phase-controllable elements across the aperture of 160 is determined, in part, by the so-called transverse coherence parameter, otherwise known as the Fried parameter, which is characteristic of the scale size of the turbulent atmosphere.\nThe spatial bandwidth of the phase modulator 160 is designed to accommodate the spatial bandwidth indicative of the wavefront distortions, 120, subject to Nyquist constraints, as is known in the art. The sampling of the wavefront sensor 170 is also designed to accommodate the wavefront distortions 120 subject to Nyquist comstraints. In the image compensation systems (to be discussed with respect to FIG. 2 below), the spatial bandwidth requirements for the corrective plate are the same, in terms of resolving the wavefront error distortions sampled by the reference beam. The imaging resolution, on the other hand, is dictated by the diffraction limit of the overall optical system. In most cases, the Fried parameter scale size of the turbulence is far greater than that of the pixel size required to faithfully image the object.\nEach of the compensation elements 150 and 160 is controlled and configured in real-time using various classes of optical detectors, algorithms and electronic networks, examples of which are feedback, feed-forward and multi-dither systems, as is known in the art. One example of an optical feedback control loop is depicted in FIG. 1. It consists of a wavefront error sensor 170, a processor module 177, and a pair of electronic drivers 180 and 185 that provide control signals to tip-tilt compensator 150 and the spatial phase modulator 160, respectively. Ideally, the driver 185 will generate a spatial phase map indicative of a wavefront-reversed replica, whose phase is given by \u2212\u03c6. The resultant beam will therefore possess a wavefront that is a combination of the incident phase distortion, \u03c6, with the correction phasemap, \u2212\u03c6, resulting in a wavefront with a net phase given as \u03c6+(\u2212\u03c6)=0, indicative of an aberration-free reference beam.\nThe optical feedback control system is designed to correct the error 140 such that the wavefront at 190 is unaberrated. This is done by driving the wavefront error seen by 170 to a minimum relative to a known reference signal that encompasses the non-common-path errors of the optical system. Upon convergence of the servo control configuration, the resultant reference beam that strikes the camera/detector 190 will be, ideally, free of wavefront errors. In this state, the overall optical receiver system 100 will provide an image of the reference beam source 110, to its diffraction limit. Given that this system functions in real-time, dynamic path distortions can be tracked and compensated, with a residual error determined by the servo-loop gain and its bandwidth.\nTurning now to FIG. 2A, a compensated image adaptive optical system 200 is shown, typical of the prior art. The goal of this system is to enable precision imaging of an extended object 205 in the presence of dynamic path distortions 220, with the resultant image viewed by camera 290. The basic adaptive optical aspect of the system functions in a manner similar to that of FIG. 1. However, in the system depicted in FIG. 2A, there are now two different input beams incident upon a telescope 230. One of the two input beams is designated as a reference beam 110, and provides the same function as that of beam 110 of FIG. 1. That is, it is in the form of a sub-diffraction-limited optical source that samples the path distortions 220. The other incident light is an image-bearing beam of object 205 whose spatial information is also distorted by the path distortions 220, and whose high-fidelity compensated image is sought.\nThe reference and image-bearing beams both traverse the same input optical components and propagation path, including the telescope 230, intermediate focal plane 235, a collimation component, represented by lens 245, tip-tilt compensator 150, spatial phase modulator 160, imaging optics 247. The reference beam 110 and the image-bearing beam 205 both impinge upon beam splitter 265. The beam splitter directs each respective input beam into a different direction. The incident reference beam 110 emerges from one port of the beam splitter as beam 266 and propagates along one direction; and, the incident image-bearing beam 205 emerges from the other port of the beam splitter as beam 267 and propagates along a second direction. The reference beam 266 is directed to the adaptive optical control loop, and the image-bearing beam 267 is directed to a camera/detector module 290. Beam splitter 265 partitions the reference and image beams using a variety of discrimination techniques including polarization, wavelength, spatial frequency, temporal gating, as is known in the art.\nIn the compensated imaging system 200, the reference beam 266 exiting beam splitter 265 is directed to an adaptive optical processor in a manner analogous to that described with respect to FIG. 1. However, as opposed to FIG. 1, in the compensated imaging system of FIG. 2A, light from the incident reference beam 110 does not strike the camera 290. The sole purpose of the reference beam in this case is to provide path-distortion information to the wavefront error sensor 270 in the servo-loop so that, upon correction of the distortions imposed posed onto the reference beam, the image-bearing beam can be viewed with little or no distortion. The feedback loop, operationally, is similar to that of FIG. 1, namely, the raw wavefront-error information is inputted into processor 175 (see 277 in FIG. 2A), which, in turn provides error correcting information to drivers 180 and 185, the outputs of which provide signals to the tip-tilt compensator and the spatial phase modulator, 150 and 160, respectively.\nThe reference beam 266 emerging from beam splitter 265 passes through an intermediate image plane 255, followed by lens 249, which transforms the beam to a pupil plane. The beam is then scaled by the telescope (lenses 247 and 249) to satisfy the spatial bandwidth constraints of the wavefront-error sensor (WFS) 270. In this system, the WFS is a so-called Shack-Hartmann class of configuration. As is known in the art the Shack-Hartmann WFS consists of a lenslet array 271 and a detector array 273, the latter positioned at the focal plane of the lenslets. This pair of elements provides a spatial mapping of the local tilt phase errors across the overall pupil-plane aperture, that characterize the path-distorted incident reference wavefront 110. As known in the art, the required number of lenslets is a function of the square of the ratio of the input aperture size to that of the coherence (Fried) parameter indicative of the incident wavefront distortions. Under these constraints, it is assumed that the incident wavefront can be described as a series of plane-wave segments, each with a different tilt, or phase slope, and all concatenated together. Hence, each plane-wave segment is considered as a diffraction-limited beamlet, each with a different tilt angle.\nFIGS. 2B and 2C, respectively, illustrate the basic prior-art principles of the Shack-Hartmann WFS, as applied to an aberrated wavefront 220, and a distortion-free wavefront 221. The WFS, identical in both FIGS. 2B and 2C, consists of a lenslet array 271 and a multi-pixel detector array 273, the latter positioned at the focal plane of the lenslets. FIG. 2B depicts the operation of the WFS assuming an input reference beam whose wavefront is aberrated. Each plane-wave segment of the input beam 222 is incident upon a different lenslet in the array 271. In the presence of no wavefront phase errors beyond the wavefront sensor's Nysquist limit, a nearly diffraction-limited sinc-squared pattern will appear at each respective focal plane. However, since each plane-wave segment is comprised of a tilted wavefront, the sinc-squared pattern at each respective focal plane at the detector array 273 win be spatially shifted, with the lateral shift increasing with the slope of the local tilt. In most systems, especially atmospheric compensation systems, the spots will not form diffraction-limited pattern due to phase errors beyond. Nyquist which distort each individual spot. A \u201cbeam's eye view\u201d at the detector surface 273, in the presence of the aberrated bean is shown in 274. Note that the array of focused spots is does not precisely overlap the grid-pattern. This is indicative of a typical aberrated beam, whose local tilts are randomly distributed. Therefore, each spot at the plane 274 has a correspondingly different offset in the (x,y) plane relative to the grid pattern. As is known in the art, the camera or ccd array 273 will require a sufficient number and density of resolvable detector pixels to measure the offset in spot position to ascertain the local tilt error with sufficient precision.\nFIG. 2C depicts the operation of the WFS assuming an input reference beam whose wavefront aberrations have been corrected. In the ideal case, the input beam 221 is a perfect plane wave, with a corresponding equi-phase surface across the entire input aperture to the WFS. As in FIG. 2B, each resolvable plane-wave segment of the input beam 223 is incident upon a different lenslet in the array 271. As before, an Airy disc pattern will appear at each respective focal plane along the detector surface 273. However, since each plane-wave segment has the same tilt (ideally, zero degrees with respect to the optical axis), each respective Airy pattern at the focal plane at the detector array 273 will be centered on its respective grid location. The \u201cbeam's eve view\u201d at the detector surface 273, in the presence of the compensated reference beam, is shown in 274. Note that the array of focused spots precisely overlaps the grid-pattern. This is indicative of an ideal plane wave, whose local tilts are identical, and a wavefront sensor with no internal aberrations. In actual implementations, there will be internal wavefront sensor error that must be measured and removed to produce reference spot/sub-image default positions. Therefore, each spot at the plane 274 has a zero spatial offset in the (x,y) plane relative to the grid pattern. It is the goal of the servo-loop adaptive optical system to drive an aberrated beam (comprised of is finite number of tilted plane-wave segments) to a converged wavefront whose differential tilts approach zero.\nIt is important to emphasize that the WFS detects only the reference beam, which, by definition, does not contain image information, other than the spatial information resulting from the intervening propagation-path distortions. Hence, based on the prior art, in order to realize an image-compensation adaptive optics system, a reference beam must be present in addition to the image-bearing beam. However, in many applications, a diffraction-limited reference beam will not always be present or practical, even in cooperative scenarios (whereby, knowledge of the existence of a reference beam or of an observer is not a drawback). And, in certain cases, a reference beam optical source may be undesirable for strategic considerations, since detection of a reference optical source by a third party can reveal the presence and/or location of a covert observer. For these and other considerations, it is desirable to realize a compensated imaging system without the need for a cooperative reference beam"}
-{"text": "This type of sensor may be remotely interrogated by connecting the input of the transducer to a radio frequency (RF) antenna. When the antenna receives an electromagnetic signal, the latter gives rise to waves on the surface of the substrate which are themselves reconverted into electromagnetic energy in the antenna. Thus, the device, consisting of a set of resonators connected to an antenna, has a response to the resonant frequencies of the resonators that it is possible to measure remotely. It is thus possible to produce remotely interrogable sensors. This possibility is a major advantage of surface acoustic waves and may be used notably in the context of pneumatic pressure sensors. This is because it is advantageous in this type of application to be able to place the sensor in the pneumatics whereas the interrogation electronics are stowed onboard the vehicle.\nAccording to the prior art, remote interrogation systems use interrogation signals in the form of pulses (typically with periods of about 25 \u03bcs) which are transmitted via an emitting antenna in the direction of a receiving antenna connected to the surface-wave sensor (referred to below, in the description, as a SAW sensor).\nA preferred frequency band for this type of system is the ISM (industrial, scientific and medical) band having a central frequency of 433.9 megahertz, the associated band width being 1.7 megahertz (MHz).\nGenerally, a remotely interrogable SAW sensor and its interrogation system may comprise, as illustrated in FIG. 1, in the simplified case of a single transducer:                an interrogation system 2; and        at least one resonator 1 comprising;                    an antenna 100; and            a comb transducer having interdigitated electrodes 11 and an SAW resonant cavity 13 characterized by its central frequency F and its quality factor Q (corresponding to the ratio between the central frequency and the width of the pass band). The cavity 13 comprises two series of reflectors that are uniformly spaced apart by a distance d. The transducer is connected to the antenna 100.                        \nThe interrogator 2 sends a long radio-frequency pulse so as to charge the resonator 1. After the emission has been stopped, the resonator discharges at its resonant eigenfrequency with a time constant \u03c4 equal to Q/\u03c0F. This discharge of the resonator forms the return echo detected by the receiver of the interrogator. Spectral analysis then allows the resonator frequency to be calculated and identified. This analysis may be carried out by algorithms based on Fourier transforms, for example an FFT (fast Fourier transform). This type of spectral-analysis treatment is particularly complex.\nA method for remotely interrogating passive SAW-type sensors has been proposed, in patent application WO 2008/015129, based on a frequency-modulation method. More precisely, this patent application discloses a method of measuring the resonant frequency of a resonator comprising the following steps:                emission, in succession, of radio-frequency signals of known carrier and modulation frequency, including the resonant frequency;        reception, via a receiving system, of response waves from the sensor; and        spectral analysis of the response waves from the sensor.        \nThe RF emission is frequency modulated with a modulation \u03c9m and an amplitude modulation typically of about a kHz.\nThe response signal of the sensor is amplitude modulated with a modulation frequency \u03c9m.\nThe measurement principle described above is based on the conversion of an (emitted) frequency modulation into an amplitude modulation via the transfer function of the resonator, as illustrated in FIG. 2. An emission modulated at an angular frequency \u03c9m (corresponding frequency fm) results in the power detector receiving a signal modulated at the frequency \u03c9m but more importantly, in a possible phase inversion of the modulated signal depending on whether it is below or above the resonant frequency.\nThe sinusoid, injected in the signal, at the frequency \u03c9m, representing the modulating signal, is either directly converted into an amplitude modulation by the transfer function of the resonator (positive slope beneath the resonant frequency, amplitude modulation in phase with the frequency modulation), or inverted (negative slope above the resonant frequency, amplitude modulation in antiphase with the frequency modulation). The intermediate point corresponds to a null contribution to the received signal at the modulation frequency fm.\nAround this frequency position, the amplitude of the component of the signal at the frequency fm varies linearly.\nThe applicant observed that when implementing an algorithm providing the function described above, namely transformation of the contribution to the modulation frequency into a signed datum around the null contribution to the modulation frequency fm, and when detecting the amplitude, although a sinusoid is emitted to generate the frequency modulation, it is possible to use only two components of the received signal, respectively at the maximum and minimum modulation frequencies, this observation allowing interrogation times to be very substantially reduced."}
-{"text": "This invention relates to a process and an apparatus for manufacturing a cast product and, more particularly, to a cast forging process and apparatus capable of manufacturing a machine component to which a high-strength is required, such as a scroll of a scroll compressor and a VTR drum without any fear of defect such as shrinkage cavities.\nFIGS. 7a and 7b are schematic sectional views illustrating a conventional indirect insertion-type squeeze casting apparatus disclosed in \"Illustrated Dictionary of Casting, Second Edition\" P. 121, edited by Japanese Association of Casting, FIG. 7a illustrating a state before casting and FIG. 7b illustrating a state immediately after casting. In these figures, reference numeral 1 is a top die, 2 is a bottom die, 2a is a molten metal reservoir disposed in the bottom die 2, 2b and 2c are a plurality of gates and cavities (a single cavity may be used), respectively, disposed between the top and the bottom dies 1 and 2, 3 is an extrusion plunger, 4 is a die-clamping ram attached to the top die 1, 5 is a knock-out punch, 6 is a molten metal and 7 is a solidified casing.\nThe operation will now be described. The top die 1 is firmly urged against and fastened to the bottom die 2 by the die-clamping ram 4. The molten metal 6 is supplied to the molten metal reservoir 2a formed in the bottom die 2 immediately before the casting. In this state, the injection plunger 3 is moved downward to pressurize the molten metal 6. The injection plunger 3 continues to pressurize the molten metal 6 at a high pressure equal to or more than 500 atmospheric pressure until the molten metal 6 solidifies and casting is completed. Then, the injection plunger 3 is moved rearward (upward) and the top die 1 is moved upward by the die-clamping ram 4 to open the die and the cast product 7 can now been taken out from the die by moving the knock-out punch 5 upward. The gate sections 2b are removed by a separate pressing step to obtain a cast product corresponding to the cavity portion 2c.\nFIGS. 8a and 8b illustrate a scroll member for use in a scroll compressor having a hollow boss portion, FIG. 8a being a front view and FIG. 8b being a sectional view taken along B--B line in FIG. 8a. In these figures, reference numeral 40 is a scroll member for a scroll compressor, which has a scroll-shaped spiral teeth portion 40a on one side of a base plate portion 40b and a boss portion 40c axis-symmetrically formed on the other side of the base plate portion 40b. This is a component which requires a high precision of the order of .mu.m.\nFIG. 9 illustrates a schematic sectional view of an indirect insertion cast forging apparatus for casting the scroll member for the scroll compressor illustrated in FIGS. 8a and 8b based on the disclosure of \"Illustrated Dictionary of Casting, Second Edition\" P. 121, edited by Japanese Association of Casting, the left half illustrating the state before casting and the right half illustrating after casting. In the figure, reference numeral 1c is a gate disposed in the top die 1 and 1d is a cavity.\nThe operation will now be described. The top die 1 is firmly urged against and fastened to the bottom die 2 by the die-clamping ram 4. The molten metal 6 is supplied to the molten metal reservoir 2a formed in the bottom die 2 immediately before the casting. In this state, the injection plunger 3 is moved upward to pressurize the molten metal 6. The molten metal 6 is introduced by the injection plunger 3 to the cavity 1d through the gate 1c. The injection plunger 3 continues to pressurize the molten metal 6 at a high pressure equal to or more than 500 atmospheric pressure until the molten metal 6 solidifies and casting is completed. Then, the die-clamping ram 4 is moved-upward to open the die 1, and the injection plunger 3 is moved upward and a connected cast product in the cavity 1d and the molten metal reservoir 2a can now been taken out from the die. The gate sections 1c are removed by a separate step of sawing to obtain a cast product 40 corresponding to the cavity portion 1d.\nAs is well known, the cast forging process is a process in which a molten metal is injected into a die cavity at a low speed capable of maintaining a laminar flow to prevent the generation of cast cavities due to the trapped gas and in which a high pressure about 500 atmospheric pressure or more is kept being applied to the molten metal until it is solidified to prevent the generation of shrinkage cavities due to the volume shrinkage upon the solidification of the molten metal, the process being known as a process for obtaining a defect-free, quality cast product. However, while it is necessary to keep applying a high pressure until each corner of the molten metal in the die solidify in order to prevent the generation of the shrinkage cavities, in the case of the indirect insertion cast forging process illustrated in FIGS. 7a and 7b, the molten metal 6 flows through the gate 2b to reach the cavity 2c and solidifies therein, so that it is necessary that the molten metal within the gate 2b does not solidify and stays in a molten state until the molten metal within the cavity 2c solidifies in order that the high pressure applied to the molten metal 6 within the molten metal reservoir 2a from the injection plunger 3 keeps being transmitted to the molten metal within the cavity 2c. Also, the molten metal within a recessed portion at the center of the cavity 2c in the illustrated example is required to solidify at a later time than the molten metal surrounding it. That is, in order to prevent the generation of the shrinkage cavity, it is necessary that a directional solidification in which the solidification is achieved progressively from the outer portion of the cavity 2c toward the molten metal reservoir 2a can be achieved. In order to achieve this, it is necessary that the gate 2b has large dimensions, which leads to problems of requiring a machine cutting of the gate 2b to obtain the cast product corresponding to the cavity 2c after the casting 7 is taken out from the die because the press cutting provides a very poor surface conditions and dimensional accuracy of the cut surface. The second problem is the limited configuration of the product, in which the dimension of the thin-wall portion for example is limited because the molten metal which becomes a product in the cavity 2c must be solidified from the outer portion toward the inner portion.\nFurther, in order to prevent generation of shrink cavities in the scroll member of the scroll compressor illustrated in FIGS. 8a and 8b, it is necessary that the directional solidification in which the solidification of the molten metal progresses from the outer portion of the cavity 1d toward the molten metal reservoir 2a can be achieved as in the apparatus illustrated in FIG. 9. In order to achieve this, firstly, it is necessary that the gate 1c has large dimensions, which leads to problems of requiring a machine cutting of the gate 1c to obtain the cast product corresponding to the cavity 1d after the casting is taken out from the die because the press cutting provides a very poor surface conditions and dimensional accuracy of the cut surface. Secondly, with this configuration of the product, the shrink cavities in the base plate portion near the gate cannot completely be eliminated even when the gate thickness is made similar to the thickness of the base plate portion."}
-{"text": "In medical practice, syringe pumps are used for drugs which need high accuracy and have a short half-life in the body. In operating theatres and intensive care units, syringe pumps are mounted in stacks of in particular six to eight pumps that, however, require a lot of space and allow only low visibility and recognition of the indication of a drug provided on the syringe. Apart from occupation of large space syringe pumps have more problems, since a nominal trumpet curve and, thus, a constancy index cannot be achieved in practice due to the fact that the plunger usually made of rubber sticks to the walls of the syringe and therefore advances in pulses rather than continuously. Further, syringe pumps also have low sensitivity in occlusion pressure reading, that becomes a problem in neonatal infusions and recently with the use of wearable bolus large volume injectors. Syringe pumps are extensively used mainly in Europe, where about 40% of the beds in each hospital are equipped with a syringe pump, and for insulin infusions and most of immunoglobulin and Parkinson's disease infusions. It has been proposed to replace insulin syringe pumps by diaphragm pumps that, however, cannot be realized in practice since insulin crystallizes and renders active and passive valves of the mechanism leak.\nPre-filled syringes are part of a growing pharmaceutical delivery sector and work well for injections, but are problematic for longer term infusions, since they become bulky especially for newer biological drugs and have a volume limit of about 60 to 100 ml. So, the pumps become bulky as the needed infusion volume increases.\nSyringe pumps are used because prior art peristaltic pumps had a low accuracy and causes a pulsatile flow and, hence, not a linear flow per infusion cycle, wherein during a part of the cycle there is no infusion but sometimes even a backflow, so that their constancy index is high for short half-life drugs. Short-term accuracy can be expressed by the concept of constancy index. This is the shortest period during a steady-state operation of a pump over which a measurement of output consistently falls within +/\u221210% of the mean rate. These data are derived from flow tests performed over 24 hours at 1 ml/h, wherein the flow is recorded at 30 seconds intervals during the final 18 hours period and the average rate is compared with the flow over each short period.\nPeristaltic pumps comprise a housing and a compressible tubing arranged within the housing. Basically, there are two different embodiments of the peristaltic pump, wherein in the one embodiment the tubing is arranged along a straight track, whereas in the other embodiment the tubing is formed as a loop resulting in a more economical design with a smaller physical size and less producing costs. The former embodiment which is called a linear peristaltic pump is mainly used nowadays, while the present invention deals with the latter embodiment. The tubing is to be filled with a fluid to be delivered from its inlet to its outlet. The fluid is caused to move through the tubing by engagement elements, typically in the form of rollers driven by rotary means such an electric motor or a mechanically driven shaft. The engagement elements cause an occlusion of the tubing by squeezing it against a wall or track within the housing so that the fluid is forced through the tubing due to the movement of the engagement elements. The use of a peristaltic pump is advantageous in so far as the fluid does not come into contact with the operating environment, which renders the peristaltic pump suitable for medical applications like infusion of drugs where it is important to avoid contact of the fluid with the environment. Further, the mechanical components of a peristaltic pump do not come into contact with the fluid. So, the components of a peristaltic pump remain free from contamination by the fluid. As a result, a peristaltic pump is easy to clean and to sterilize because the tubing can be simply discarded after use, and a new tubing can be provided for the next use.\nHowever, a disadvantage is that it is difficult with a peristaltic pump to achieve a constant or pulseless flow of the fluid through the tubing. Pulses are created when the engagement elements disengage from the tubing and, therefore, the occlusion is removed with the result of that a void is generated in the disengagement region of the tubing. Namely, in this region the tubing returns to its normal round shape resulting in an increase of the volume which are filled by fluid from the outlet of the tubing. This leads to a reduction of the flow rate of the fluid at the outlet of the tubing for the duration of the pulse.\nIn other words, the pulsatile behavior of a rotary peristaltic pump results from the alternation of a forerunner or leading engagement element by the next follower or trailing engagement element. V is the volume which is encapsulated between two neighboring engagement elements. In case of a rotary peristaltic pump, with \u03c6 defining an angular position so that the angular position of the trailing engagement element is \u03c61 and the angular position of the leading engagement element is \u03c62, the volume V encapsulated between both these engagement elements extends along an angular distance which is defined by the difference between both the aforementioned angular positions \u03c62 and \u03c61, i.e. \u0394\u03c6=\u03c62\u2212\u03c61. +\u0394V/\u0394\u03c6 represents an increase of the volume defining a so-called frontwave which is displaced in front of each engagement element and is advanced by it. \u2212\u0394V/\u0394\u03c6 represents a decrease of the volume defining a so-called depression which arises behind each engagement element. The enclosed volume V between two neighboring engagement elements squeezing the tubing with unchanged distance between them is constant so that the fluid is just transported, if +\u0394V/\u0394\u03c6 and \u2212\u0394V/\u0394\u03c6 are the same so thatV+\u0394V/\u0394\u03c6\u2212\u0394V/\u0394\u03c6=V, resulting in that the pressure remains constant as well.\nIf otherwise an upstream portion of the tubing is larger than a downstream portion so that it is\u0394V2/\u0394\u03c6>\u0394V1/\u0394\u03c6 andV+\u0394V2/\u0394\u03c6\u2212\u0394V1/\u0394\u03c6=V+Vdifference,wherein \u0394V1/\u0394\u03c6 represents an increase of volume defining a front wave in front of the forerunner or leading engagement element and \u0394V2/\u0394\u03c6 represents an increase of volume defining a front wave in front of the next follower or trailing engagement element, the pressure is increased by elastic tubing forces of the inflated portion due to increase of volume.\nAt the moment the leading engagement element stops squeezing the tubing and a front/back communication through a thin film of fluid is established under it, the frontwave +\u0394V/\u0394\u03c6 suddenly disappears (so that it cannot push fluid anymore) and is replaced by the frontwave in front of the trailing engagement element which now takes over relay of infusion. Also \u2212\u0394V/\u0394\u03c6 disappears behind it. But due to the disengagement of the leading engagement element a new additional difference volume \u0394Vd/\u0394\u03c6 is built up and continues to be present until the disengagement of the leading engagement element is fully completed. A void creating the aforementioned new additional volume difference \u0394Vd arises, as the resilient tubing becomes asymptotic or starts with a larger diameter disengagement curvature at this point, resulting in a creation of a negative pulse \u0394V/\u0394\u03c6\u2212\u0394Vd/\u0394\u03c6 in the flow graph (V,\u03c6) (where \u0394Vd depends on the geometry of the engagement element and the disengagement curvature \u0394r/\u0394\u03c6 of the housing and r is the radius of the disengagement curvature increasing so as to let the engagement element disengage).\nThe flow rate is defined as \u0394V/\u0394t. By multiplying the nominator and the denominator each with \u0394\u03c6, the above equation can be written asflow rate=(\u0394V/\u0394\u03c6)\u00b7(\u0394\u03c6/\u0394t).\nThis makes evident that for a constant rotational movement \u0394\u03c6/\u0394t, if \u0394V is constant at the vicinity of an engaging element at any location along its movement path, the flow rate will be constant without any variation or pulse.\nThe lack of a constant fluid caused by pulses in the tubing renders a peristaltic pump unsuitable for certain precision applications. E.g., in applications where a small volume of fluid is required, such as where less than a complete revolution of the rotor is used, the effect of the pulses are particularly disadvantageous.\nU.S. Pat. No. 5,533,878 A discloses a squeeze type pump wherein the resilient tubing has a larger diameter at the start of an infusion cycle than at its end.\nU.S. Pat. No. 7,654,127 B2 discloses a peristaltic pump with increased dimensions in an upstream portion of the tubing so as to increase pressure inside the tubing before disengagement of a roller squeezing the tubing. However, the pressure is suddenly released when the leading roller is going to be disengaged from the tubing. Therefore, a perturbation of flow happens first with a sudden increase of flow and second with a decrease of it at same quantity due to disengagement of the roller from the tubing after some rotational degrees.\nU.S. Pat. No. 3,826,593 A proposes the provision of a cam compressing the tubing at another point at the same time the engagement roller is disengaged from the tubing in a peristaltic pump. However, this solution requires higher costs and a larger number of assembly parts and is therefore not suitable for a disposable pumping mechanism.\nU.S. Pat. No. 5,470,211 A teaches a roller pump with the provision of a controlled curvature at the inlet and the outlet of the tubing.\nIt is an object of the present invention to provide a continuously operating pulseless and accurate peristaltic pump having a small size so that it can replace many bulky syringe pumps in limited space."}
-{"text": "Tertiary amides are used as solvents in a variety of industrial chemical processes, and their recycling for reuse in these processes is of significant commercial value. Depending on the type of reaction performed in the tertiary amide, such recycling may require the removal of reaction byproducts not easily separated from it in a cost-effective manner. One example of such byproducts includes carboxylic acids. In some cases, these must be removed from used solvent as part of recycling, and in other cases they must be removed from the solvent while it is still in use to dissolve process intermediates. The latter situation frequently compounds the difficulty of removing the acids, since potentially sensitive intermediates may be present. The commercial preparation of sucralose is an example of a process where carboxylic acids need to be removed from a tertiary amide solvent, both alone and in the presence of sensitive process intermediates.\nSucralose (4,1\u2032,6\u2032-trichloro-4,1\u2032,6\u2032-trideoxygalactosucrose), a high-intensity sweetener made from sucrose, can be used in many food and beverage applications.\nA number of different synthesis routes for the preparation of sucralose have been developed in which the reactive hydroxyl in the 6 position is first blocked with an acyl group to form a sucrose-6-ester. The acyl group is typically acetyl or benzoyl, although others may be used. The sucrose-6-ester is then chlorinated to replace the hydroxyls at the 4, 1\u2032 and 6\u2032 positions to produce 4,1\u2032,6\u2032-trichloro-4,1\u2032,6\u2032-trideoxygalactosucrose 6-ester (referred to herein as sucralose-6-ester), followed by hydrolysis to remove the acyl substituent and thereby produce sucralose. Several synthesis routes for formation of the sucrose-6-ester involve tin-mediated acylation reactions, with illustrative examples being disclosed in U.S. Pat. Nos. 4,950,746; 5,023,329; 5,089,608; 5,034,551; and 5,470,969, all of which are incorporated herein by reference.\nVarious chlorinating agents may be used to chlorinate the sucrose-6-ester, and most commonly a Vilsmeier-type salt such as Arnold's Reagent will be used. One suitable chlorination process is disclosed by Walkup et al. (U.S. Pat. No. 4,980,463), incorporated herein by reference. This process uses a tertiary amide, typically N,N-dimethyl formamide (\u201cDMF\u201d), as the chlorination reaction solvent. After the chlorination is complete, adducts of Arnold's reagent on the base sucrose moiety and excess chlorinating reagent are neutralized (\u201cquenched\u201d) with aqueous base to provide the sucralose-6-ester in an aqueous solution, accompanied by the tertiary amide solvent and salts resulting from reactions of the chlorination reagent. The sucralose-6-ester is then deacylated to produce sucralose. One suitable deacylation process is taught by Navia et al, U.S. Pat. No. 5,498,709, the entire disclosure of which is incorporated herein by reference.\nIn such processes, carboxylic acids need to be removed at various points in the process. Accordingly, facile means of removing these acids are of commercial value."}
-{"text": "1. Field of the Invention\nThis invention relates to apparatus for supporting a conduit from a beam and more particularly to a beam clamp coated with a moisture, resistant polymeric material in which a clamping member is secured to a U-bolt that supports a conduit adjacent to the beam with polymeric encapsulated nuts engaging the threads of the U-bolt is sealed relation to prevent corrosion of the threaded connection.\n2. Description of the Prior Art\nBeam clamps are well known in the art of securing a conduit, pipe, or the like in a preselected position to a structural member, such as an I-beam. A variety of types of beam clamps are known to position the conduit either vertically, parallel, or on edge relative to the longitudinally extending horizontal flange of the I-beam. U.S. Pat. No. 2,338,006 is an example of a clamp device releasably engageable with the vertical edge of a supporting structure to position the conduit in spaced parallel relation to the supporting structure. This type of clamp employs a threaded member for securing a U-shaped strap by a clip to the wall of the supporting structure. The clip is secured to the supporting structure by the threaded connection of a bolt and nut. The nut engages the threaded portion of the bolt that extends through the supporting structure.\nOne of the disadvantages of this type of arrangement is that threaded connections are subject to corrosive damage when employed in a corrosive environment. In many cases a beam clamp is located in a corrosive environment where moisture can easily enter the threaded connection of nuts and bolts. The corrosive damage can \"freeze\" the threaded connection. Consequently the beam clamp can become rendered useless by creating diffucult maintenance problems in disengaging the nut from the bolt to disassemble the clamp from the beam and the conduit.\nIt is also well known in the art as disclosed in U.S. Pat. Nos. 3,724,706; 3,784,236; and 3,799,584 to coat components of an assembly that must be capable of efficient assembly and disassembly in a corrosive environment with a moisture resistant, insulating, resilient, polymeric material, such as a polyvinylchloride (PVC). By providing male and female coupling members with a moisture, resistant polymeric coating a moisture resistant seal is formed around the otherwise exposed threads of the coupling members. In this manner moisture is prevented from entering the point of engagement, for example, of a nut on a threaded bolt to prevent freezing of the nut on the bolt.\nWhile it is known to coat the components of a beam clamp, i.e. the clamping member and the U-bolt, with a moisture resistant, polymeric coating, the threaded connection of the clamping member to the U-bolt is exposed. The components generally are assembled in the field where the nuts and bolts are exposed to the affects of moisture.\nTo resist the deleterious affects of moisture it is known to apply a moisture resistant, polymeric material by brush coating the exposed nuts and the U-bolt. However, a moisture resistant coating applied in the manner to the exposed threaded connections and fasteners in a corrosive environment has proved unsatisfactory in preventing corrosion. Brushing on a polymeric coating after installation will not deter the accummulated affects of the corrosive atmosphere prior to the application of the coating. Furthermore, if improper attention is given to the brushing application of polymeric material to the exposed threads and nuts, an insufficient coating thickness to prevent corrosion may result.\nTherefore, there is need to provide for a beam clamp a connection of the clamping member and the conduit support member capable of preventing corrosive damage while permitting rapid assembly and disassembly of the beam clamp components."}
-{"text": "1. Field of the Invention\nThis invention relates generally to a balancer for dehydration tubs for use in washing machines or the like, and more particularly to such a balancer of the liquid-in-container type including a concentric multi-compartment container.\n2. Description of the Prior Art\nIn dehydration tubs used in washing machines having a centrifugal dehydrating function or the like, a liquid-in-container type balancer is mounted on an upper open end of the dehydration tub for the purpose of correcting an unbalanced condition thereof due to one-sided laundry. For improvement in a correcting force, the interior of the balancer is divided by one or more vertically extending concentric partition walls into a plurality of compartments. A predetermined amount of liquid (usually, salt water) is contained in each of the compartments.\nFIGS. 23 and 24 illustrate the construction of a conventional balancer of the type described above. Referring to FIG. 24, the balancer comprises an annular balancer container 1 having in its interior a partition wall 2 standing from the bottom thereof. The interior of the container 1 is divided by the partition wall 2 into inner and outer compartments 3 and 4. A lid 5 is attached to the top of the container 1 so as to close upper openings of the compartments 3 and 4. The lid 5 has two inlets 6 and 7 formed therein to correspond to the compartments 3 and 4 respectively. A liquid 8 is poured through the inlets 6 and 7 into the respective compartments 3 and 4. Thereafter, closures 9 and 10 are closely fitted into the inlets 6 and 7 respectively.\nIn the balancer described above, the liquid 8 contained in the container 1 leaks out or flows between the compartments 3 and 4 if the compartments are not watertightly sealed by the lid 5 or if the inlets 6 and 7 are not watertightly closed by the respective closures 9 and 10. Consequently, an expected correcting force cannot be obtained when the liquid 8 leaks out of the container 1 or flows between the compartments 3 and 4. In view of this problem, watertight tests need to be carried out for the balancer.\nThe watertight tests are carried out in the following procedure. The interior of the compartment 3 is pressurized or depressurized through the inlet 6 after the lid 5 has been attached to the container 1. Consequently, the watertightness is tested at a boundary 11 between an inner wall of the container 1 and the lid 5 and at a boundary 12 between the partition wall 2 and the lid 5. Subsequently, the interior of the compartment 4 is pressurized or depressurized through, the inlet 7 in order that the watertightness is tested at a boundary 13 between an outer wall of the container 1 and the lid 5. Finally, the inlets 6 and 7 are closed by the respective closures 9 and 10 after the liquid 8 has been poured into the compartments 3 and 4 through the inlets respectively. The completed balancer is then put into a chamber, and the interior of the chamber is pressurized or depressurized in order that the watertightness at boundaries between the circumferential edges of the inlets 6 and 7 and the closures 9 and 10 is tested.\nIn the conventional balancer, the watertight test needs to be carried out for each of the compartments 3 and 4, and furthermore, the other watertight test needs to be carried out at a stage of the end product of balancer. Accordingly, steps of the watertight test is increased with an increase in the number of compartments of the container. Furthermore, since the portions of the balancer to be tested are diverse, an equipment for the watertight test is rendered complicate, and accordingly, the cost of equipment is increased.\nOn the other hand, the number of compartments of the container or the height of the balancer is increased in the prior art so that the performance of the balancer is improved. In each case, the size of the balancer is increased, which results in an increase in the size of the washing machine or a decrease in a washing capacity of the washing machine.\nFurthermore, the correcting performance of the balancer mainly depends upon the configuration of the container 1. The correcting performance of the balancer cannot be altered after its configuration has been determined."}
-{"text": "1. Field of the Invention\nThe present invention relates to an active energy ray-curable composition and an antistatic film.\n2. Description of the Related Art\nGenerally, resin materials exhibit excellent electrical insulation properties and are therefore very useful for applications requiring electrical insulating properties such as insulators, but surfaces thereof are easily charged with static electricity, which is thus likely to result in dust adsorption and electrostatic troubles.\nTo solve such problems, it is known to use various antistatic agents (in particular, cationic antistatic agents).\nFor example, JP2011-225797A discloses an \u201cantistatic coating composition including an antistatic copolymer (A) having a quaternary ammonium base in the molecule, a compound (B) having three or more ethylenically unsaturated groups, a photopolymerization initiator (C) containing a photopolymerization initiator (c-1) having a solubility in water of 0.2 g/L or less and a photopolymerization initiator (c-2) having a solubility in water of 1 g/L or more, and a solvent (D) containing an alcohol-based solvent (d-1)\u201d ([claim 1]), and also discloses that the photopolymerization initiator (C) is \u201cat least one selected from the group consisting of 1-hydroxy-cyclohexyl-phenyl-ketone, 2-methyl-1-[4-(methylthio)phenyl]-2-morpholinopropan-1-one, bis(2,4,6-trimethylbenzoyl)-phenylphosphine oxide, 2-hydroxy-1-{4-[4-(2-hydroxy-2-methyl-propionyl)-benzyl]-phenyl}-2-methyl-propan-1-one, 2,4,6-trimethylbenzoyl-diphenyl-phosphine oxide, oligo {2-hydroxy-2-methyl-1-[4-(1-methylvinyl)phenyl]propanone}, and 2,4,6-trimethylbenzophenone\u201d ([claim 4])."}
-{"text": "A plunger lift is an apparatus that is used to increase the productivity of oil and gas wells. In the early stages of a well's life, liquid loading is usually not a problem.\nWhen rates are high, the well liquids are carried out of the tubing by the high velocity gas. As the well declines, a critical velocity is reached below which the heavier liquids do not make it to the surface and start to fall back to the bottom exerting back pressure on the formation, thus loading up the well. A plunger system is a method of unloading gas in high ratio oil wells without interrupting production. In operation, the plunger travels to the bottom of the well where the loading fluid is picked up by the plunger and is brought to the surface removing all liquids in the tubing. The plunger also keeps the tubing free of paraffin, salt or scale build-up. A plunger lift system works by cycling a well open and closed. During the open time a plunger interfaces between a liquid slug and gas. The gas below the plunger will push the plunger and liquid to the surface. This removal of the liquid from the tubing bore allows an additional volume of gas to flow from a producing well. A plunger lift requires sufficient gas presence within the well to be functional in driving the system. Oil wells making no gas are thus not plunger lift candidates.\nAs the flow rate and pressures decline in a well, lifting efficiency declines geometrically. Before long the well begins to \u201cload up\u201d. This is a condition whereby the gas being produced by the formation can no longer carry the liquid being produced to the surface. There are two reasons this occurs. First, as liquid comes in contact with the wall of the production string of tubing, friction occurs. The velocity of the liquid is slowed, and some of the liquid adheres to the tubing wall, creating a film of liquid on the tubing wall. This liquid does not reach the surface. Secondly, as the flow velocity continues to slow the gas phase can no longer support liquid in either slug form or droplet form. This liquid along with the liquid film on the sides of the tubing begin to fall back to the bottom of the well. In a very aggravated situation, there will be liquid in the bottom of the well with only a small amount of gas being produced at the surface. The produced gas must bubble through the liquid at the bottom of the well and then flow to the surface. Because of the low velocity very little liquid, if any, is carried to the surface by the gas. Thus, as explained previously, a plunger lift will act to remove the accumulated liquid.\nA typical installation plunger lift system 100 can be seen in FIG. 1 (prior art). Lubricator assembly 10 is one of the most important components of plunger system 100. Lubricator assembly 10 includes cap 1, integral top bumper spring 2, striking pad 3, and extracting rod 4. Extracting rod 4 may or may not be employed depending on the plunger type. Below lubricator 10 is plunger auto catching device 5 and plunger sensing device 6. Sensing device 6 sends a signal to surface controller 15 upon plunger 200 arrival at the well top. Plunger 200 is shown to represent the plunger of the present invention and will be described below in more detail. Sensing the plunger is used as a programming input to achieve the desired well production, flow times and wellhead operating pressures. Master valve 7 should be sized correctly for tubing 9 and plunger 200. An incorrectly sized master valve will not allow plunger 200 to pass. Master valve 7 should incorporate a full bore opening equal to the tubing 9 size. An oversized valve will allow gas to bypass the plunger causing it to stall in the valve. If the plunger is to be used in a well with relatively high formation pressures, care must be taken to balance tubing 9 size with the casing 8 size. The bottom of a well is typically equipped with a seating nipple/tubing stop 12. Spring standing valve/bottom hole bumper assembly 11 is located near the tubing bottom. The bumper spring is located above the standing valve and can be manufactured as an integral part of the standing valve or as a separate component of the plunger system.\nSurface control equipment usually consists of motor valve(s) 14, sensors 6, pressure recorders 16, etc., and electronic controller 15 which opens and closes the well at the surface. Well flow \u2018F\u2019 proceeds downstream when surface controller 15 opens well head flow valves. Controllers operate on time, or pressure, to open or close the surface valves based on operator-determined requirements for production. Modern electronic controllers incorporate features that are user friendly, easy to program, addressing the shortcomings of mechanical controllers and early electronic controllers. Additional features include battery life extension through solar panel recharging, computer memory program retention in the event of battery failure, and built-in lightning protection. For complex operating conditions, controllers can be purchased that have multiple valve capability to fully automate the production process.\nIn these and other wells it is desirable to measure the downhole temperature and pressure versus time, chemical profiles and other data. This information is used to figure oil and gas reserves and production plans. Conventional methods include dropping special sensors called pressure bombs via cable down the tubing. Pressure bombs can be attached to the wireline or left downhole to be retrieved by fishing at a later date. Special trucks with a crew are used which is expensive for the well operator.\nIn FIG. 2 (prior art), a special truck called a wireline (also called slickline) rig 200 is used to drop a downhole equipment data logger (temperature and/or pressure and/or time) 207 down tubing 266 of the well. Nominally the tubing is two inches in diameter, and data logger 207 is about three feet long. Wireline rig 200 has an on-board computer 201 for data recording. Hoistable crane 202 supports electric line 206 which usually requires a lubricator 203 and a blowout protector 204. A spool and hoist assembly 205 controls electric line 206. All this special equipment is costly to lease for the well operator. Furthermore, the use of this equipment requires the complete shutdown of the well during the operation of dropping special data logger 207.\nWhat is needed is an improved data logger sensor that can be dropped down a well and retrieved without a wireline rig. The plunger will house and deliver the data logger to the bottom of the well to take readings. Then the well operator can turn the well on to flow the plunger and data logger to the surface without the use of a wireline rig and crew. This sensor should be easily detachable to the plunger and readily plugged into a computer to retrieve the measured downhole temperature and/or pressure. The present invention fulfills these needs for the well operator/producer."}
-{"text": "An infusion pump assembly may be used to infuse a fluid (e.g., a medication or nutrient) into a user. The fluid may be infused intravenously (i.e., into a vein), subcutaneously (i.e., into the skin), arterially (i.e., into an artery), and epidurally (i.e., into the epidural space).\nInfusion pump assemblies may administer fluids in ways that would be impractically expensive/unreliable if performed manually by nursing staff. For example, an infusion pump assembly may repeatedly administer small quantities of an infusible fluid (e.g., 0.1 mL per hour), while allowing the user to request one-time larger \u201cbolus\u201d doses."}
-{"text": "1. Field of the Invention\nThe present invention relates to a network game in which a plurality of players participate to progress with the game, and, more particularly, to execution of an event in which a plurality of players participate.\n2. Description of the Related Art\nThe recent advancement on the network technology has made network games (online games) vivid in each of which a plurality of players participate to progress with the game. A system that executes a network game includes a server apparatus which is managed by one who runs the network game, and client devices of individual players which are connectable to the server apparatus over a network such as the Internet.\nAn RPG (Role Playing Game) among the network games progresses as multiple players participate in the game at the same time so that their player characters cooperate with one another to battle with an opponent character, or the players make their player characters battle with one another. Recently, in particular, an MMORPG (Massively Multiplayer Online RPG) which permits participation of more players is becoming popular. A network game such as an MMORPG requires that the movements of the player characters of all the participating players should be transmitted to the client devices manipulated by the players without contradiction.\nJapanese Patent Application Laid-Open Publication No. 2008-99906 discloses a technique of allowing a server apparatus to centralize management on the timings for the actions and processing of the player characters, and periodically synchronize with the timing of progressing the game. The technique disclosed in Japanese Patent Application Laid-Open Publication No. 2008-99906 prevents the individual client devices from processing time-sequential information with some delays which would bring about inconsistency on the progress of the game. This scheme can be adapted to processing of an event in which a plurality of players participate in a network game.\nWhen an event occurs in a network game, a movie may be played back on the individual client devices as a previous step to the event so that the individual players watch the movie at the same time. In a network game such as an MMORPG in which multiple players participate, however, there are a variety of players including skilled or experienced players and beginners. In a case where experienced players and beginners participate in the same event, is likely that the experienced players have already seen such a movie so that they may feel tedious about being tied on watching the movie for a long time whereas the beginners are apt to see it for the first time.\nFrom the viewpoint of just playing back a movie, the each client device may be permitted to freely skip the movie as each player is likely to skip some scenes in progressing with a stand-alone game, so that the player is relieved of a long tie-up. When an event such as a battle immediately follows the end of a movie, however, the timing for a process related to the event in the progress of the game may vary among the individual client devices. From the viewpoint of playing an event such as a battle with a plurality of players in cooperation, it is not desirable that while some players are watching a movie, the other players freely start the battle."}
-{"text": "1. Field of the Invention\nThis invention relates to scroll fluid machine, in which sucked fluid is compressed with stationary and revolving scrolls and discharged to the outside.\n2. Description of the Related Art\nA scroll fluid machine compresses fluid sucked from its peripheral part in a sealed space formed by its stationary and revolving scrolls progressively as the fluid is fed toward its central part, and discharges the compressed fluid from the central part. As the fluid is compressed, the temperature in the sealed space formed by the wraps is elevated. This poses a problem that bearings, seal members, etc. provided in drive parts are soon deteriorated. Heretofore, the scrolls are cooled to hold the temperature within a predetermined temperature.\nWell-known cooling systems cool either a non-driven part, i.e., the stationary scroll, or a driven part, i.e., the revolving scroll.\nFIG. 16 shows a technique concerning a non-driven part cooling system. As shown, a revolving scroll 116 which is mounted on a frame 109 provided in a sealed housing 105, comprises a disc-like body 114 having a shaft 113 depending therefrom. The frame 109 has a central hole, in which a drive shaft 104 coupled to a drive (not shown) is fitted for rotation, and the shaft 113 is eccentrically coupled to the drive shaft 104. The revolving scroll 116 has a wrap 115 engaging with a wrap 111 of a stationary scroll 112.\nThe stationary scroll 112 has a peripheral wall having a suction hole 118. When the revolving scroll 116 is revolved relative to the stationary scroll 112 with the rotation of the drive shaft 104, a sealed space formed by the wraps 111 and 115 is progressively reduced in volume, thus compressing gas entering the sealed space. The compressed gas is discharged from a discharge hole 121 formed in a central part of the stationary scroll 112 through a discharge pipe 120 to the outside.\nA plurality of radially spaced-apart heat pipes 122 are provided in the body 110 of the stationary scroll 112 to remove heat generated in a compression stroke as described above.\nFIG. 17 shows a well-known cooling system for cooling driven part, i.e., the revolving scroll.\nA housing 211 as shown comprises a rear and a front housing part 212 and 213, and a drive shaft 214 is supported for rotation by bearings 215 in a bearing portion of the rear housing part 212. The drive shaft 214 has an extension projecting outward from the bearing portion and coupled to a motor (not shown). The drive shaft 214 also has an eccentric portion 214b, which has an eccentric axis 02--02 with respect to the axis 01--01 of the drive shaft 214 by a distance .delta..\nA revolving scroll 216 which is coupled to the eccentric portion 214b of the drive shaft 214, has a disc-like plate 216a having a mirror finished front surface, a spiral wrap 216b formed on the front side of the mirror finished plate 216a, a boss 216c formed as the driving center with an axial line 02--02 on the rear side of the plate 216a and having smaller diameter than the inner peripheral surface edge of above portion 213b, a ring-like ridge 216d formed on the rear side of the above 216a and on the periphery thereof, and a plurality of radial vent holes 216e formed in a diameter direction above the ridge 216d.\nA stationary scroll 221, which is secured to the front housing part 213, has a disc-like plate 211a having a mirror finished rear surface, a spiral wrap 221b formed on the rear side of the plate 211a and a peripheral wall 221c surrounding the wrap 221b.\nThe wraps 216b and 221b of the revolving and stationary scrolls 216 and 221 engage with or wrap each other at a predetermined deviation angle, and they form a plurality of compression chambers or spaces when the revolving scroll 216 is revolved.\nThe drive shaft 214 has a counterweight 225 mounted on its portion extending in the rear housing part 212, and a centrifugal fan 226 is mounted on the counterweight 225 to generate cooling air flow with the rotation of the drive shaft 214.\nIn the prior art non-driven part cooling system shown in FIG. 16, in which the heat pipes 122 are provided in the stationary scroll body, the heat absorbing portions of the heat pipes 122 are more remote from the revolving scroll which is driven than from the stationary scroll. Therefore, the neighborhood of the bearings, seal members and other parts which are driven in contact with the revolving scroll 116 in the driving thereof, is cooled less efficiently compared to the cooling of the stationary scroll. This means that uniform temperature distribution cannot be obtained.\nThe heat radiating portions of the heat pipes 122 are cooled by their heat radiation to the sealed housing inner space 105a, which is filled with gas-sucked through a suction pipe 119.\nIn communication with the space 105a is the suction hole 118, through which gas enters the compression space which is formed by the stationary and revolving scrolls. This means that gas having been elevated in temperature by the heat radiation from the heat pipes 122 again enters the compression space through the suction hole 118, thus reducing the cooling efficiency.\nIn order to prevent the cooling efficiency reduction, it is necessary to provide special cooling means on an external part to which the suction pipe 119 is connected, thus complicating the construction and increasing the size of the apparatus.\nIn the well-known driven part cooling system shown in FIG. 17, with the rotation of the drive shaft 214 external gas is sucked through a suction passage 227 by the centrifugal fan 226 and led through a ring-like space B and a cooling air passage 220 to be discharged through a discharge passage 228.\nSince in this system the gas having cooled down a central part of the revolving scroll 216 is discharged along the rear side of the revolving scroll 216 and through the discharge passage 228, the provision of the discharge passage is necessary. In addition, in order to increase the cooling efficiency, a cooling fan for cooling the rear side of the stationary scroll 221 has to be provided, thus increasing the size of the apparatus."}
-{"text": "Adhesively backed index tabs are known in the art. The tabs typically have a coating of removable or permanent pressure sensitive adhesive on the back, such that the tab may be placed on the edge of the page or other substrate. The tabs may be used, for example, to mark a page in a book. A tab may have a portion that extends beyond the edge of the page, to identify the page that is marked. The tab may have a writable portion, such that the user can write or print text and/or graphics.\nIndex tabs of this type are typically supplied on a sheet, from which the user may remove the index tabs as desired. The sheets are typically U.S. letter size or similar and are not always convenient to store due to the relatively large size. Alternatively the tabs may be supplied affixed to sheets of a pad. A sheet of tabs that is removed from a pad is often easily misplaced because it is thin. Alternatively, the tabs may be supplied in a dispenser, with the user pulling the tabs from the dispenser as needed. The dispensers are typically small and easily misplaced or somewhat bulky such that they cannot be conveniently kept where most needed."}
-{"text": "1. Field of the Invention\nThe present invention relates to a gas sensor of Severinghaus type using a hydrogen ion sensitive FET transducer of the structure of an oblong gate-insulated field effect transistor having its gate part at the front end and its electrode part at the other end.\n2. Description of the Prior Art\nThe measurement of concentration of gas such as carbon dioxide or ammonia gas is, of course, important in industrial applications, but recently in the medical field, the measurement of partial pressure of gas in living bodies begins to be taken seriously. For example, in medicine, continuous measurements of gas partial pressure in blood of anesthetized patients, those with advanced diseases or those in the convalescent stage are lending themselves to discovery of emergent situations. For such purposes, a very small gas sensor with less than 2 mm diameter which may be inserted in any blood vessel or muscle tissue is required.\nFor the aforementioned purpose, heretofore in use has been a gas sensor of the Severinghaus type using a very small glass electrode. Minuaturization of the glass electrode, however, is known to involve following problems:\n(a) Because the resistance of the glass film is approx. 10 M.OMEGA., an amplifier with a high input impedance is necessary. PA1 (b) The glass film, being thin, has low mechanical strength. PA1 (c) Because of the electrode area being small, the resistance of the glass film is large. PA1 The measuring instruments should be large and complex and the electrodes themselves are brittle and tend to break. Therefore, especially as a sensor to be inserted into tissues of living bodies for measuring gas partial pressure in living bodies, this instrument has posed a problem in practical applications. PA1 (a) Because of the use of the solid electrode, sensor is inflexible. PA1 (b) Its electric resistance is increased due to the miniaturization. PA1 (1) The photosensitivity of pH-ISFET exists not only on its gate part but on whole of the Si substrate including its electrode part. Accordingly, it is difficult to totally eliminate the photosensitivity by coating only the gate part with a black polymer, as above described. PA1 (2) In the Severinghaus gas sensor, the partial pressure of gas is measured by utilizing the minute change of pH of the gas absorbing liquid. Accordingly, it is necessary to avoid changes of pH of the gas absorbing liquid due to factors other than the gas partial pressure. However most of such black pigments, being chemically unstable, will become acidic due to autoxidation, etc., during a long period of storage, often inviting changes in characteristics as a gas sensor. Accordingly, using the pH-ISFET with its gate part coated with a hydrophilic polymer colored black for allaying the photosensitivity of the pH-ISFET is problematical in practical applications as a gas sensor.\nOn the other hand, a carbon dioxide gas sensor using a solid pH electrode of metal oxide in place of the glass electrode is disclosed in U.S. Pat. No. 3,719,576, etc. This sensor is smaller and slenderer than that using a glass electrode and is, therefore, suitable as a sensor to be inserted in living tissues, but has the following disadvantages:\nOn this ground, the miniaturization had its limit.\nThis problem has been solved by making use of a hydrogen ion sensitive FET transducer (hereinafter referred to as pH-ISFET) having the gate-insulated field effect transistor structure described in U.S. Pat. No. 4,218,298 in place of the glass electrode of the solid electrode. Such a gas sensor using the pH-ISFET is, as described in U.S. Pat. No. 4,409,980, a gas sensor composed of a pH-ISFET and a reference electrode deposited on the substrate in proximity to the gate part of the pH-ISFET, an insulator tube housing this pH-ISFET and the reference electrode with lead wires connected thereto, the gate part of the aforementioned pH-ISFET being located at an opening part provided in the insulator tube, and the lead wires extended along the tube, an electrical insulation resin closing the tube by filling the space of the part inside the tube wall housing the lead wire connecting parts, a hydrophilic polymer layer containing electrolytes which undergo change in hydrogen ion concentration, as it absorbs the gas, and which is placed around the gate part of the pH-ISFET and the reference electrode, enveloping both of them, and a gas permeable membrane coating at least whole of the aforementioned polymer layer.\nThe aforementioned gas sensor is preferable as a gas sensor to be inserted in living bodies. However, the pH-ISFET usually has a photo-sensitivity of the order of 10.sup.-5 -10.sup.-6 V/lux. Therefore, when a gas sensor using the pH-ISFET is used for monitoring, etc., during operation, when it receives illumination of several thousands luxes - several tens of thousands luxes, the photosensitivity of the pH-ISFET poses a serious obstacle to its practical use.\nCoating the gate part of the pH-ISFET with a hydrophilic polymer colored black for allaying the photosensitivity of the pH-ISFET is described in U.S. Pat. No. 4, 273,636. However, the aforementioned proposal is to allay the photosensitivity of ion sensor using ISFET and does not relate to a gas sensor. It is in principle practicable to allay the photosensitivity of a gas sensor by utilizing a pH-ISFET with a hydrophilic polymer colored black coated on its gate part. In this instance, the gate part of the pH-ISFET and the reference electrode are enclosed by a hydrophilic polymer containing a gas absorbing liquid colored black and the outside of this polymer is coated with a gas permeable membrane. It has become evident, however, that diminishing the photosensitivity of a gas sensor by this method involves the following problems:\nThe present inventors have found out the improved gas sensor of this invention as a result of investigations carried out for providing a practically useful gas sensor in which the photosensitivity of the gas sensor of the Severinghaus type using pH-ISFET is allayed."}
-{"text": "When disasters (e.g., hurricane, storm, tsunami, etc.) strike, many establishments set up collection mechanisms to collect donations to support the victims. These mechanisms may include links on web sites to encourage donations electronically. For example, donations may be done via credit cards or via any other forms of monetary transfers including transfers from accounts associated with payment facilitators (e.g., PayPal Inc. of San Jose, Calif.). Unfortunately, when there are people making donations to the victims, there are also scammers that take advantage of the situations. These scammers may fraudulently claim that they collect the donations for the victims. Since the donation levels typically are at their highest immediately after the disasters struck, it is often difficult to quickly verify the scammers or where the donations go during that time. By the time the scams are detected, large amounts of donations have already been misdirected to the scammers and away from the victims."}
-{"text": "1. Field of the Invention\nThis invention relates to control systems and, more particularly, to methods of driving dynamic and steady-state behavior of a process toward more optimum operating conditions.\n2. Discussion of Related Art\nMultivariable predictive control algorithms, such as DMCplus\u2122 from AspenTech\u00ae of Aspen Technologies, Inc. or RMPCT from Honeywell International Inc., are a combination of calculations to drive the dynamic and steady-state behavior of a process toward a more optimum operating condition. The steady-state algorithm used in the control scheme is most commonly a linear program (LP), but sometimes is a quadratic program (QP). For small problems, understanding the LP or QP solution is relatively simple. Two-dimensional problems can be visualized on a paper and demonstrated to an operator to gain understanding of the process. With some detailed modeling background, engineers and well trained operators can understand medium-sized problems (less than 10 dimensions). However, larger, more interactive problems, often require offline simulation. This can take a significant amount of time to understand, even qualitatively.\nTypically, an operator of a multivariable predictive controller (MPC) can observe current constraints and may have access to an open loop model of the process. However, to fully understand the constraint set relief, the operator would need a detailed understanding of the process model and the ability to trace independent and dependent relationships through the model. For that, an offline simulation or analysis tool is required. Otherwise, the operator cannot know how much to change a constraint or which constraint is the next to become active.\nOne concept for an offline simulation uses a matrix pivot in which unconstrained manipulated variables (MVs) are swapped with constrained controlled variables (CV). The constraints become \u201cindependents,\u201d and the unconstrained variables become \u201cdependents.\u201d The matrix pivot can be symbolized as follows:\n      [                                        y            1                                                            y            2                                ]    =                              [                                                    A                                            B                                                                    C                                            D                                              ]                \u2061                  [                                                                      x                  1                                                                                                      x                  2                                                              ]                    \u2062                          [                                                  x              1                                                                          y              2                                          ]        =                  [                                                            A                                  -                  1                                                                                                      -                                      A                                          -                      1                                                                      \u2062                B                                                                                        CA                                  -                  1                                                                                    D                -                                                      CA                                          -                      1                                                        \u2062                  B                                                                    ]            \u2061              [                                                            y                1                                                                                        x                2                                                    ]            \nHowever, this approach does not provide quantitative answers as to how much any operator change will affect the controller solution.\nThere is a need for a simple utility that can analyze past or current dynamic matrix control (DMC) solutions of any-sized problem, in real-time, to provide an operator meaningful, quantitative instructions for DMC controller constraint relief."}
-{"text": "Integrated circuits are formed on semiconductor substrates, or wafers. The wafers are then sawed into microelectronic dies (or \u201cdice\u201d), or semiconductor chips, with each die carrying, for example, a respective integrated circuit (e.g., a microprocessor) or a microelectromechanical system (MEMS) device (e.g., an accelerometer). In some examples, each semiconductor chip is mounted to a package or carrier substrate using either wirebonding or \u201cflip-chip\u201d connections. The packaged chip is then typically mounted to a circuit board, or motherboard, before being installed in a system, such as an electronic or a computing system.\nWhile lower frequency applications typically use direct conductive connections (e.g., vias and solder formations) to connect to the circuit board, higher frequency applications, such as millimeter wave devices, often use electromagnetic coupling and a waveguide structure. In order to achieve the desired performance in the higher frequency applications, the package substrate must be precisely manufactured, which typically involves manufacturing processes that are not compatible with conventional semiconductor manufacturing processes. As a result, manufacturing costs are increased.\nAccordingly, it is desirable to provide an improved method for forming a microelectronic assembly for use in high frequency applications that allows the required performance to be achieved while minimizing costs. Other desirable features and characteristics of the invention will become apparent from the subsequent detailed description and the appended claims, taken in conjunction with the accompanying drawings and the foregoing technical field and background."}
-{"text": "1. Field of the Invention\nThe present invention relates to magnetic disk drives, and more particularly to a mechanism for preventing damaging contact between a magnetic head slider assembly and a rigid magnetic disk due to shock loads applied to a magnetic disk drive.\n2. History of the Prior Art\nAmong the better known data storage devices are magnetic disk drives of the type in which a magnetic head slider assembly floats on an air bearing at the surface of a rotating rigid magnetic disk. Such disk drive are often of the so-called Winchester type in which one or more rigid magnetic disks are located within a sealed chamber together with one or more magnetic head slider assemblies. The magnetic disk drive may include one, two or more rigid magnetic disks, and the slider assemblies may be positioned at one or both of the opposite sides of each of the magnetic disks.\nEach magnetic head slider assembly in magnetic disk drives of the type referred to is typically coupled to the free outer end of a different one of a plurality of elongated arms or load beams. The slider assembly is mounted in a manner to permit gimballed movement at a free outer end of the arm so that an air bearing between the slider assembly and the surface of the rigid magnetic disk can be established and maintained. The elongated arm is coupled to an appropriate mechanism for moving the arm across the surface of the disk so that a magnetic head contained within the slider assembly can address different ones of concentric data tracks on the disk for writing information into or reading information from the data tracks.\nAn example of an arm assembly having a gimballed mount for a magnetic head slider assembly is provided by U.S. Pat. No. 3,931,641 of Watrous. The arm assembly described in the Watrous patent includes a relatively rigid load beam having a rigid bearing member at a free outer end thereof for receiving a protuberance on a spring element. The spring element is spot welded to the load beam and has an end thereof defining a flexure. The flexure includes a pair of stiff crosslegs mounted on an opposite pair of flexible outer fingers and a central finger. The central finger mounts a magnetic head slider assembly, and gimballed movement is provided by the load protuberance on the spring element which is held in contact with the bearing member at the end of the rigid load beam. Such arrangement provides desired gimballing action by allowing a pitch and roll of the slider assembly around mutually orthogonal axes while at the same time resisting radial, circumferential and yaw motions.\nIn magnetic disk drives of the type described, physical contact of the surface of the magnetic disk containing the data tracks by the slider assembly must be avoided. Consequently arrangements must be provided for preventing such contact when the magnetic disk is at rest or is otherwise not rotating at its nominal operating speed. When the disk is rotating at its nominal operating speed, the air bearing between the slider assembly and the disk is usually sufficient to prevent contact therebetween. One technique commonly employed is to move the slider assembly onto a portion of the disk where physical contact can be tolerated whenever the disk is decelerated to rest. One or more so-called parking zones are provided on the surface of the disk where the slider assembly may rest. When the disk is accelerated to its nominal operating speed, the air bearing again forms and the slider assembly may then be moved to the data tracks. Another common technique for preventing contact between the slider assembly and the disk surface when the disk is decelerated to rest is to pivot the arm so as to move the slider assembly away from the surface of the disk. This moves the slider assembly into an unloaded position well away from the disk surface. When the disk is again accelerated to the nominal operating speed, the arm is pivoted so as to bring the slider assembly back into a loaded position in which the air bearing forms between the slider assembly and the surface of the disk.\nAn example of an arrangement for raising and lowering the arm so as to move the slider assembly between loaded and unloaded positions is provided by a copending application of Warrent L. Dalziel, Ser. No. 759,900, filed July 29, 1985 and commonly assigned with the present application. The Dalziel application describes a mechanism for loading and unloading one or more slider assemblies using simple mechanical apparatus. Such apparatus engages the arms to raise them into unloaded positions and at the same time limit their radial movement. A separator element is movable between first and second positions to selectively engage the arms. When the separator element is in a first position, the slider assemblies are in loaded positions, and radial motion of the slider assemblies is not inhibited. Movement of the separator element into a second position moves the slider assemblies into unloaded positions and at the same time locks an associated carriage assembly in a selected position to inhibit radial motion of the slider assemblies.\nEach time the disks are decelerated to rest in the arrangement described in Dalziel application Ser. No. 759,900, the slider assemblies are raised into the unloaded positions and at the same time radial movement of the arms is limited. This does much to protect the sensitive components of the magnetic disk drive against damage due to shock loads such as may occur if the disk drive is dropped or otherwise bumped. Nevertheless, depending upon the manner of coupling the slider assemblies to the ends of the elongated arms, it may be difficult or impossible to prevent damaging contact between the slider assemblies and the disk surfaces. This is so even though the arms are pivoted to move the slider assemblies into unloaded positions well away from the disk surfaces. Inadvertent dropping of the magnetic disk driven onto a hard surface or other application of severe shock loads may still result in the slider assemblies contacting the disk surfaces with possible damage to the slider assemblies or the disk surfaces or both. Typically the slider assembly is coupled to the free outer end of the arm by a flexure element which is highly flexible and which therefore permits substantial movement of the slider assembly away from the arm in response to shock loads and other substantial forces.\nAccordingly, it is an object of the invention to provide an improved magnetic disk drive, and particularly an improved disk drive of the type in which one or more magnetic head slider assemblies are employed in conjunction with one or more rigid magnetic disks.\nIt is a further and more specific object of the invention to provide an improved magnetic disk drive in which the possibility of damaging contact of the disk surfaces by the slider assemblies in response to shock loads when the arm is raised into an unloaded position is greatly minimized or eliminated."}
-{"text": "Enterprise content management (ECM) covers a broad range of applications, including document management (DM), Web content management (WCM), records management (RM), digital asset management (DAM), search of managed content, and the like. A content management system (CMS) suitable for managing the various content (also referred to herein in some examples as \u201cfiles\u201d or \u201cdocuments\u201d) that an enterprise produces or generates, retains or otherwise stores, manipulates or modifies, etc. can support the requirements of one or more of such applications, and optionally other requirements, to provide a coherent solution in which content processes, management processes, and the like are capable of accessing content across a variety of applications subject to access controls, permissions, and the like. Content managed by a CMS can include one or more of documents, images, photos, Web pages, records, XML documents, other unstructured or semi-structured files, etc. Content retained in a CMS can also include directory structures such as folders, file trees, file plans, or the like, which can provide organization for multiple content items in addition to storing or otherwise representing relationships between content item, etc. An \u201centerprise\u201d can generally refer to an organization, such as for example a business or company, a foundation, a university, or the like, and can have content requirements related to one or more business processes, content uses, etc.\nA CMS manages the actual digital binary content, the metadata that describes a context of each content item, associations between a content item and other content or content items, a place and classification of a content item in a repository, indexes for finding and accessing content items, etc. The CMS can also manage processes and lifecycles of content items to ensure that this information is correct. The CMS can also manage one or more workflows for capturing, storing, and distributing content, as well as the lifecycle for how long content will be retained and what happens after that retention period.\nA CMS for use in enterprise content management can include one or more of document management tools, applications, and interfaces to support general office work, search, and discovery. Workflow management capabilities of a CMS can support numerous business processes, optionally including, but not limited to, case management and review and approval. Collaboration applications and services of a CMS can support the collaborative development of information and knowledge in the creation and refinement of content and documents. This collaborative development of information and knowledge can be achieved through providing access to content managed by the CMS to multiple users. To prevent conflicting or discontinuous editing streams, a user can be allowed to check out or lock content for modification and check in the modified content such that other users are prevented from editing content concurrently. Web content management services of a CMS, which can be scalable, can support the delivery and deployment of content from the enterprise to its customers. Records management capabilities of a CMS can capture and preserve records based upon government-approved or other standards. A standards-based platform can also provide access to applications that use these standards, such as publishing, image management, email management, etc."}
-{"text": "1. Technical Field\nThe present invention relates to a wireless communication device.\n2. Background Technology\nAs a method for establishing a network, there is a method for performing Wi-Fi Direct (hereinafter referred to as \u201cWi-Fi Direct Connection\u201d) under Wi-Fi Alliance standard development. A wireless communication device that is possible to perform the Wi-Fi Direct Connection can directly communicate with other wireless communication devices that are also possible to perform the Wi-Fi Direct Connection as one-on-one communication or a multiple devices communication without using an access point. Also, it is possible to communicate with a wireless communication device that is possible to perform a wireless LAN connection method used with Wi-Fi (hereinafter referred to as \u201cWi-Fi Connection\u201d) function as an access point. In addition, since the Wi-Fi Direct Connection uses the Wi-Fi technologies, the high-speed communication is possible. By using the Wi-Fi Direct Connection, for example, it is possible to directly send an image to a printer from a camera or a cellular phone and print it out.\nJapanese Laid-open Patent Publication No. 2011-166417 (Patent Document 1) and http://205.149.128.22/Wi-Fi_Direct.php (Non-patent Document 1) are examples of the related art."}
-{"text": "Bednorz and Muller, Z. Phys. B64, 189 (1986), disclose a superconducting phase in the La-Ba-Cu-O system with a superconducting transition temperature of about 35 K. The presence of this phase was subsequently confirmed by a number of investigators [see, for example, Rao and Ganguly, Current Science, 56, 47 (1987), Chu et al., Science, 235, 567 (1987)]. Chu et al., Phys. Rev. Lett., 58, 405 (1987), Cava et al., Phys. Rev. Lett., 58, 408 (1987), Bednorz et al., Europhys. Lett., 3, 379 (1987)]. The superconducting phase has been identified as the composition La.sub.l-x (Ba,Sr,Ca).sub.x O.sub.4-y with the tetragonal K.sub.2 NiF.sub.4 -type structure and with x typically about 0.15 and y indicating oxygen vacancies.\nWu et al., Phys. Rev. Lett., 58, 908 (1987), disclose a superconducting phase in the Y-Ba-Cu-O system with a superconducting transition temperature of about 90 K. The compounds investigated were prepared with nominal compositions (Y.sub.1-x Ba.sub.x).sub.2 CuO.sub.4-y and x=0.4 by a solid-state reaction of appropriate amounts of Y.sub.2 O.sub.3, BaCO.sub.3 and CuO in a manner similar to that described in Chu et al., Phys. Rev. Lett., 58, 405 (1987). This reaction method comprised heating the oxides in a reduced oxygen atmosphere of 2.times.10.sup.-5 bars (2 Pa) at 900.degree. C. for 6 hours. The reacted mixture was pulverized and the heating step was repeated. The thoroughly reacted mixture was then pressed into 3/16 inch (0.5 cm) diameter cylinders for final sintering at 925.degree. C. for 24 hours in the same reduced oxygen atmosphere.\nHundreds of other papers have since disclosed similar solid state reaction processes. Other papers have disclosed various solution and precipitation methods for preparing the reactants to be heated at temperatures of 800.degree.-850.degree. C. and above.\nHirano et al., Chemistry Letters, 665, (1988), disclose a process for producing Y-Ba-Cu-O superconductors by the partial hydrolysis of a solution of Ba metal, Y(O-iPr).sub.3 and Cu-acetylacetonate or Cu-alkoxides in 2-methoxy or 2-ethoxy ethanol. The solution was stirred in dry nitrogen and heated at 60.degree. C. for 12 hours. The solution was then hydrolyzed by the slow addition of water diluted with solvent. Stirring and heating continued for several hours. Stirring continued while the solution was evaporated under vacuum at about 60.degree. C. and an amorphous precursor powder was obtained. The powder was calcined in flowing oxygen at temperatures between 800.degree. and 950.degree. C. for up to 24 hours. The calcined powder was pressed and sintered in flowing oxygen at temperatures up to 920.degree. C. and then annealed at temperatures between 450.degree. and 550.degree. C.\nIt is highly desirable to form precursors that can be used to produce powders that have small size particles, i.e., generally sub-micron in size, and that can be pressed into desired shapes, sintered and converted to superconducting Y-Ba-Cu-O."}
-{"text": "Provisioning a computing system can be a complicated and error-prone task. The computing system may need to be provisioned with many applications, each of which may need to have settings configured, data entered, and/or files altered. Provisioning multiple computing systems across an entire data center or even several data centers may be significantly more complicated and may involve the tedious and repetitive task of configuring the same enterprise-wide settings again and again. To reduce the complexity and rate of error, many provisioning systems now enable administrators to provision new systems from stored backup images of a fully configured system. Provisioning a new system from a backup image may allow administrators to quickly provision new systems both locally and remotely with a much reduced chance of configuration errors.\nUnfortunately, traditional systems for provisioning new systems from backup images may suffer from a variety of inefficiencies. Some traditional systems may aggregate read requests to the backup system in order to improve efficiency. However, sub-optimal read request batch sizes may lead to slow responses from the backup system, increasing the amount of time it takes to provision a new system from a backup and leading to user frustration. Accordingly, the instant disclosure identifies and addresses a need for additional and improved systems and methods for dynamically adjusting batch request sizes."}
-{"text": "In the design, processing and implementation of embedded capacitors for power decoupling in electronic appliances and systems, different techniques have been proposed for the implementation of the embedded capacitors. In the implementation of embedded capacitors, limitations may arise in relation to the integration of the embedded capacitors with power distribution networks in terms of performance and space utilization. For example, blind/buried vias are typically required in existing techniques, resulting in process limitations and increased manufacturing complexity. Furthermore, in existing techniques, vias for high current and high-frequency decoupling, respectively, must typically be provided, resulting in increased number of interconnection and solder joints.\nOften, the different existing techniques for the implementation of embedded capacitors are designed and limited to a specific process technology during manufacturing, which can impose further limitations of existing techniques due to lack of adaptability for different manufacturing process technology."}
-{"text": "1. Field of the Invention\nThis invention relates to building modules for the erection of walls, more particularly for freestanding, sound absorbing/reducing walls and retaining walls. These are particularly modular elements comprising at least two subelements, preferably in the form of beams arranged at an angle to one another. The modules may be regarded also as a building module kit from which at least one or more complete modules can be assembled. Accordingly, the term \"building module\" is intended to comprise unitary modules as well as multi-pan and more complex modules.\nIn many cases, one of the subelements or beams is constructed as a longitudinal beam extending substantially parallel to the wall plane and comprising at least two cross-beams arranged at an angle to one another. The cross beams extend between a pair of the longitudinal beams, in a direction transverse to the wall plane. Two of said cross-beams form support blocks with planar top and bottom sides which act as spacing means between the individual modules. Each module is designed to be held in place by gravity, its own weight, and the weight of a bulk filling material. Each module normally stacks upon another module to build up wall height, and is installed adjacent to another in an end-to-end arrangement to build wall length.\n2. Background Art\nThe typical retaining wall or freestanding wall is usually constructed with a plurality of flame-like elements comprising at least two subelements, preferably in the form of beams arranged at an angle to one another and connected in a form locking or material locking manner. U.S. Pat. No. 4,384,810, to Neumann, describes a locking beam that forms a three dimensional lattice in a construction system for plantable shoring walls. Said lattice comprises support blocks with planar top and bottom sides which are stacked one above the other. The blocks act as spacing means between individual planting levels. The structure includes a longitudinal component consisting of a base plate and a breast part. The longitudinal component and the locking beam extend parallel to the wall plane and always rest on two support blocks. The support blocks are spaced apart a distance and extend in a direction transverse to the wall plane. Each of the consecutive stacks of the support block pairs, together with the longitudinal component, the locking beam, and the earth within, form a construction section acting as a static slope shoring unit.\nU.S. Pat. No. 5,017,050, to Jaecklin describes building elements for supporting a grid wall with a bulk filling material comprising at least two subelements, preferably in the form of beams arranged at an angle to one another, and connected in a form locking or material locking manner.\nThere is at least one hole or recess in the cross beam which is open over a part of its circumferential contour and adapted to receive a longitudinal beam in order to establish a form locking connection between the beams. The beams are secured against separation and displacement from one another in a mounted state. The contour of the hole or recess overlaps the other side of the longitudinal beam so as to form an abutment. The recess is shaped so as to permit a partial lateral insertion of the longitudinal beam while establishing contact between the abutment and the longitudinal beam in a first rotational position. The longitudinal beam is then rotated downward through an arc until it is fully inserted into the recess. Once fully inserted, a form locking connection is established which cannot be reversed when the construction is under load.\nU.S. Pat. No. 5,181,351 to Jaecklin, describes building elements for supporting grid walls with a bulk material filling. The frame-like elements of the invention comprise at least two subelements, preferably in the form of beams arranged at an angle to one another, and connected in a form locking or material locking manner.\nOne of the subelements or beams is constructed as a longitudinal beam extending substantially parallel to the wall plane and comprises at least two profile legs arranged at an angle to one another. The first of these profile legs forms at least one bearing surface for the bulk filling material, while the second of these profile legs forms a retaining surface for the bulk filling material facing the inside space of the frame.\nSuch structural elements or structural systems have heretofore not been concerned with the need for structural features which allow for the adequate supply and retention of water necessary to sustain maximum plant growth without costly replanting. It is to this need and to other deficiencies of the prior art that this invention is directed."}
-{"text": "Personal audio devices, including wireless telephones, such as mobile/cellular telephones, cordless telephones, mp3 players, and other consumer audio devices, are in widespread use. Such personal audio devices may include circuitry for driving a pair of headphones or one or more speakers. Such circuitry often includes a speaker driver including a power amplifier for driving an audio output signal to headphones or speakers."}
-{"text": "This invention relates generally to gas turbine engines and more particularly to variable stator vane assemblies for use in such engines.\nGas turbine engines operate by combusting a fuel source in compressed air to create heated gases with increased pressure and density. The heated gases are ultimately forced through an exhaust nozzle, which is used to step up the velocity of the exiting gases and in-turn produce thrust for driving an aircraft. The heated air is also used to drive a turbine for rotating a fan to provide air to a compressor section of the gas turbine engine. Additionally, the heated gases are used for driving rotor blades inside the compressor section, which provides the compressed air used during combustion. The compressor section of a gas turbine engine typically comprises a series of rotor blade and stator vane stages. At each stage, rotating blades push air past the stationary vanes. Each rotor/stator stage increases the pressure and density of the air. Stators serve two purposes: they convert the kinetic energy of the air into pressure, and they redirect the trajectory of the air coming off the rotors for flow into the next compressor stage.\nThe speed range of an aircraft powered by a gas turbine engine is directly related to the level of air pressure generated in the compressor section. For different aircraft speeds, the velocity of the airflow through the gas turbine engine varies. Thus, the incidence of the air onto rotor blades of subsequent compressor stages differs at different aircraft speeds. One way of achieving more efficient performance of the gas turbine engine over the entire speed range, especially at high speed/high pressure ranges, is to use variable stator vanes which can optimize the incidence of the airflow onto subsequent compressor stage rotors.\nVariable stator vanes are typically circumferentially arranged between an outer diameter fan case and an inner diameter vane shroud. A synchronizing mechanism simultaneously rotates the individual stator vanes in response to an external actuation source.\nIn some situations, it is advantageous to divide the compressor section into upper and lower halves to expedite maintenance of the gas turbine engine. It is particularly advantageous, for example, in military applications when maintenance must be performed in remote locations where complete disassembly is imprudent. However, in dividing the compressor section into halves, the synchronizing mechanism must also be split apart. This creates two synchronizing mechanisms that must be actuated in unison to orchestrate simultaneous operation of all of the stator vanes. Synchronizing mechanisms that are located on the outer case can be accessed and spliced together easily. However, this is not the case for inner diameter synchronizing mechanisms, which cannot be accessed after assembly to attach the synchronizing mechanisms together. Thus, there is a need for an apparatus for coordinating actuation of split inner diameter synchronizing mechanisms."}
-{"text": "Quality of sleep is essential to good human health and well-being. Many types of beds have been developed, including mattresses and box springs, water beds, air mattresses, and memory foam mattresses, e.g. Mattresses and box springs comprising spring systems with coverings of pliable foam layers and possibly other materials have been in use for many years and are widely available. Many approaches and enhancements to beds and bedding have also been developed including mattress top coverings that include specially formed materials and foam to improve comfort. A wide range of materials, coverings, and configurations of blankets, sheets, pillows, and other bedding accessories are available.\nIn spite of a wide diversity of beds that are available, serious problems exist for substantially all types and configurations. Mattresses including spring systems sag over time and bugs (especially bed bugs), mites, human secretions, and other infestations and contaminants build up. Hence, many bed systems are not hygienic because there is no convenient way to launder, clean, or otherwise restore them. While a mattress with heavy coverings may be warm and cozy in the winter months, a lighter mattress covering that allows air to pass may be far more comfortable in the heat of summer. Conventional spring mattresses are also large and bulky, making them hard to store and move.\nWater beds are heavy, create damage if a leak or rupture occurs, and some persons have difficulty sleeping in them due to the motion of the water. Water beds normally include a heating system to warm the water in them that normally takes several hours to establish a requested temperature. Hence, while the water bed temperature is usually adjustable, a person may find a given temperature not comfortable and would need to wait several hours for the water bed to adjust. For at least the foregoing reasons, water beds are generally impractical for use in hotels.\nAir mattresses are normally not preferred for bedding as they don't allow air to pass, resulting in a warm and sweaty sleeping environment. In addition, some persons are unable to find comfortable sleeping positions on air mattresses.\nSome mattresses include viscoelastic foam. Viscoelasticity provides both viscous and elastic properties that enable a mattress to conform and support a person sleeping on it. The viscous nature of the material is sensitive to temperature and is normally configured so that the warmth of a person's body causes the viscoelastic foam to yield and conform; however, air does not easily pass through, making the mattress too warm for some persons, and viscoelastic mattresses are also susceptible to contaminations and infestations. Further, viscoelastic foam mattresses are generally heavy and also expensive relative to a conventional spring mattress.\nA key problem with all existing types of mattresses and beds is that while they may be tailored to persons of specific heights, weights, body shapes, and other personal aspects; it is impossible for them to be optimized for all persons. Hence, a mattress that provides stiffer springs in the region of a person's back for better back support of larger and heavier persons cannot help but generate an uncomfortably too stiff feeling for other persons who are smaller and lighter. Adjustable beds may provide some accommodations for certain sleeping positions, but if a person changes position the bed may need to be re-adjusted. Accordingly, adjustable beds substantially retain the issues with the various types of mattresses that they may be used with, and also tend to be heavy and expensive to produce."}
-{"text": "The present invention is directed to the field of hydraulic dampers for reducing undesired pressure surges or oscillations in systems requiring hydraulic damping. This invention has particular applicability to the damping of periodically varying motions in systems which result in noise or oscillation in hydraulic pressure.\nIn a particular example, undesirable oscillation in a hydraulic system occurs in the nose wheel assembly of an aircraft, which is typically steered using a hydraulic system. The pivot point of the wheel is typically located at a point ahead of the wheel axle, rather like a caster. The mass of wheel is behind the pivot, and the dynamics of the system are such that the wheel tends to \"shimmy\" while moving along the ground, in a manner similar to other commonly-observed caster-type wheels, e.g. the wheels on a shopping cart.\nSuch wheel shimmy results in wear and tear on the wheel and supporting strut which may increase exponentially resulting structural damage to the aircraft. Suppression of shimmy often requires compromises in steering control accuracy. It is desirable to provide a hydraulic means of damping which does not compromise steering accuracy.\nOne known hydraulic damper uses a piston and a pair of springs to form a compliance, and a hydraulic orifice to form a damping control which reduces rapid shimmy motion while permitting the relatively slow motion of steering. However, such dampers include a number of mechanical elements.\nHydraulic systems are prone to pressure surges which occur when large valves are quickly opened or closed. This phenomenon, also known as \"water hammer,\" is also observed in common household plumbing when faucets are abruptly turned off, producing a sound and a shock to the pipes. Such surges also create additional wear and tear in hydraulic systems, thereby shortening the operational life of the respective components."}
-{"text": "The present invention relates generally to electronic packaging. More particularly, the present invention relates to a marked structure such as a wafer or an array of packages.\nAs is well known to those of skill in the art, integrated circuits, i.e., electronic components, are fabricated in an array on a wafer. The wafer is then cut, sometimes called diced, to singulate the integrated circuits from one another.\nFIG. 1 is a cross-sectional view of a section of a wafer 10 being cut from a front-side surface 10F of wafer 10 in accordance with the prior art. Formed in wafer 10 were integrated circuits 12. Integrated circuits 12 were delineated by scribe lines 14, which included a first scribe line 14A and a second scribe line 14B, on front side surface 10F of wafer 10. For example, scribe lines 14 were formed by selective etching of a silicon oxide layer 18 on front-side surface 10F.\nTo illustrate, first scribe line 14A delineated a first integrated circuit 12A from a second integrated circuit 12B. Each scribe line 14 had a width WF.\nA back-side surface 10B of wafer 10 was attached to a tape 20. Wafer 10 was then sawed with a saw blade 22. Saw blade 22 was aligned with scribe lines 14 using an optical alignment system in a well-known manner. Saw blade 22 cut through wafer 10 along scribe lines 14. In this manner, integrated circuits 12 were singulated. Tape 20 supported wafer 10 during sawing and supported the singulated integrated circuits 12 after sawing was complete.\nGenerally, width WF of scribe lines 14 was sufficient to accommodate the width of the saw cut plus tolerance in the positioning of saw blade 22. Stated another way, width WF of scribe lines 14 was sufficiently large such that the saw cut made by saw blade 22 was always within a scribe line 14. For example, saw blade 22 is within scribe line 14B in FIG. 1.\nSince the optical alignment system used scribe lines 14 directly to align saw blade 22, saw blade 22 was aligned to scribe lines 14 to within tight tolerance. Accordingly, scribe lines 14 were relatively narrow and, more particularly, were only slightly wider than saw blade 22. To illustrate, width WF was within the range of 0.002 inches (0.051 mm) to 0.008 inches (0.203 mm).\nIn certain instances, it was important to protect the front-side surface of the wafer during sawing, e.g., from shards and particulates generated during sawing. To protect the front-side surface, the wafer was sawed from the back-side surface of the wafer as discussed below in reference to FIG. 2.\nFIG. 2 is a cross-sectional view of a section of a wafer 30 being cut from a back-side surface 30B of wafer 30 in accordance with the prior art. To protect a front-side surface 30F of wafer 30, front-side surface 30F was attached to a tape 32. Tape 32 supported wafer 30 during sawing.\nSaw blade 22 was aligned with scribe lines 14-1 on front-side surface 30F of wafer 30 using a two-step process. First, tape 32 was aligned with scribe lines 14-1. Front-side surface 30F was attached to tape 32. Tape 32 had area greater than the area of front-side surface 30F such that tape 32 had an exposed region, which extended beyond wafer 30. Tape 32 had alignment marks in the exposed region of tape 32. As an example, see alignment holes 30a and 30b of Roberts, Jr. et al., U.S. Pat. No. 5,362,681, which is herein incorporated by reference in its entirety. In the above manner, scribe lines 14-1 were aligned with the alignment marks of tape 32.\nSecond, saw blade 22 was aligned with the alignment marks of tape 32. Wafer 30 was then sawed with saw blade 22 from back-side surface 30B. However, since saw blade 22 was aligned indirectly to scribe lines 14-1 using alignment marks of tape 32, a large tolerance was associated with the alignment of saw blade 22 to scribe lines 14-1.\nTo accommodate this large tolerance, each of scribe lines 14-1 had a relatively large width WB. More particularly, referring now to FIGS. 1 and 2 together, width WB of scribe lines 14-1 of wafer 30, which was designed to be cut from back-side surface 30B, was significantly larger than width WF of scribe lines 14 of wafer 10, which was designed to be cut from front-side surface 10F. To illustrate, width WB was approximately 0.012 inches (0.305 mm) or more.\nDisadvantageously, forming scribe lines 14-1 with a relatively large width WB resulted in less integrated circuits 12 for any given size wafer 30 than the corresponding number of integrated circuits 12 formed in the same size wafer 10, i.e., there was a loss of yield of integrated circuits 12 from wafer 30. As a result, the cost of each integrated circuit 12 from wafer 30 was greater than the cost of each integrated circuit 12 from wafer 10. However, it is desirable to minimize the cost of each integrated circuit 12.\nIn accordance with the present invention, a method includes identifying and determining a position of a scribe grid on a front-side surface of a wafer with a camera. Based on this information, a computer aims a laser at a first location on a back-side surface of the wafer. The laser is fired to form a first alignment mark on the back-side surface of the wafer. Advantageously, the alignment mark is positioned with respect to the scribe grid to within tight tolerance, e.g., to within 0.001 inches (0.025 millimeters) or less.\nThe front-side surface of the wafer is attached to a tape to protect the front-side surface of the wafer during sawing. A saw blade is aligned with a scribe line of the scribe grid using the alignment mark on the backside surface of the wafer. The wafer is cut from the back-side surface along the scribe line with the saw blade.\nAdvantageously, the wafer is cut from the back-side surface thus protecting the front-side surface of the wafer and, more particularly, the integrated circuits. Of further importance, the saw blade is precisely aligned to the scribe line using the alignment mark such that the scribe line is not fabricated with the extra large width of scribe lines of conventional wafers designed to be cut from the back-side surface.\nRecall that in the prior art, in certain instances, it was important to cut the wafer from the back-side surface. However, to accommodate the large tolerance associated with back-side wafer cutting, the wafer designed to be cut from the back-side surface was formed with relatively wide scribe lines. Disadvantageously, forming the scribe lines with a relatively large width resulted in less integrated circuits for any given size wafer, i.e., a loss of yield. This resulted in a substantial increase in the cost of the integrated circuits.\nIn stark contrast, the wafer is cut from the back-side surface in accordance with the present invention without the associated loss of yield of the prior art. As a result, the integrated circuits of the wafer are protected during singulation yet are fabricated without the associated substantial increase in cost of the prior art.\nIn accordance with another embodiment of the present invention, an array of packages is marked. In accordance with this embodiment, a back-side surface of the array is scanned by a camera to identify and determine the position of fiducials on the back-side surface. Based on this information, a computer aims a laser at a first location on a front-side surface of the array. The laser is fired to form an alignment mark on the front-side surface of the array. Advantageously, the alignment mark is positioned with respect to the fiducials to within tight tolerance, e.g., to within 0.001 inches (0.025 millimeters) or less.\nThe back-side surface of the array is attached to a tape. A saw blade is aligned with the array using the alignment mark as a reference. The array is cut with the saw blade thus singulating the packages.\nA pick and place machine removes the packages from the tape. Advantageously, the packages are directly removed from the tape by a standard and relatively simple pick and place machine. Accordingly, removal of the packages from the tape is relatively simple and thus low cost. As a result, the packages are fabricated at a low cost.\nIn the prior art, an array of packages was singulated from the back-side surface. More particularly, the array was placed upside down on the tape such that a layer of encapsulant of the array was adhered to the tape and the fiducials extended upwards and were exposed. The array was singulated by cutting from the back-side surface using the fiducials as a reference.\nHowever, after singulation, the singulated packages had to be removed from the tape and inverted, e.g., using a pick and place machine with flip capability. The singulated packages had to be removed from the tape and inverted so that the singulated packages could be loaded into the grid carrier with the contacts (or other interconnection structure) facing downwards into the grid carrier. Back-end processing, e.g., automated attachment to the printed circuit mother board or automated testing, required that the singulated packages be loaded into the grid carrier in this manner. Disadvantageously, removing the singulated packages from the tape and inverting the packages required complex machinery, was labor intensive and, accordingly, increased the cost of the packages. In contrast, removal of the packages from the tape in accordance with the present invention is relatively simple and thus low cost.\nIn accordance with one particular embodiment, a structure includes a substrate, e.g., a wafer, having a first surface and a second surface. The structure further includes a reference feature, e.g., a scribe grid, on the first surface and at least one alignment mark on the second surface. The alignment mark has a positional relationship to the reference feature.\nIn accordance with another embodiment of the present invention, an array of packages includes a substrate having a first section, an electronic component such as an integrated circuit attached to the first section, metallizations on a first surface of the first section, and contacts on the metallizations. Bond pads of the integrated circuit are electrically connected to the contacts by bond wires. A layer of encapsulant covers the integrated circuit, the bond pads, the bond wires, the contacts, and the metallizations. A fiducial is on a first surface of the array and an alignment mark is on a second surface of the array. More particularly, the alignment mark is in the layer of encapsulant.\nIn one embodiment, a method includes identifying a reference feature on a first surface of a substrate such as a wafer or an array of packages. The method further includes marking a first location on a second surface of the substrate with a first alignment mark. The first alignment mark is used to determine a position of the reference feature.\nThese and other features and advantages of the present invention will be more readily apparent from the detailed description set forth below taken in conjunction with the accompanying drawings."}
-{"text": "This disclosure relates generally to the field of optics and, more specifically, to a method and apparatus for synthesizing ultra-wide bandwidth optical waveforms.\nAs waveforms are synthesized with increasing larger bandwidths, radio frequency (RF) design tends to grow in complexity. RF component design criteria such as power, amplitude ripple, gain flatness, phase distortion over bandwidth become increasingly more difficult to maintain with increasing bandwidths. What is need is an improved mechanism to increase waveform bandwidth of a radio frequency waveform during double sideband suppressed carrier modulation."}
-{"text": "The present invention relates to semiconductor storage devices, and in particular, to a semiconductor storage device having a nonvolatile memory for executing a collective erase verify operation.\nNonvolatile semiconductor storage devices have been remarkably developed in recent years, and in particular, the rewriting times of flash memories have been substantially shortened. However, about a half of the rewriting time is consumed by the verify operation of checking whether or not the data in a memory cell has been rewritten in conformity to an expected value, and therefore, it is indispensable to provide a device for reducing this verify operation time.\nIn view of the above, there has been an attempt at reducing the verify operation time by providing sense amplifiers corresponding in number to memory cells to be subjected to a simultaneous write operation and simultaneously verifying the memory cells subjected to the simultaneous write operation in the verify operation.\nAccording to the above verify method, in regard to the write operation of the above rewrite operation, the number of memory cells to be subjected to the simultaneous write operation correspond to about 1 bit to 4 kbits. Since the data to be written into individual memory cells differ from one another, it can be considered efficient to simultaneously verify the memory cells of about 1 bit to 4 kbits subjected to the simultaneous write operation.\nHowever, in regard to the erase operation of the rewrite operation, memory cells of not less than 2 kbits to 512 kbits are collectively erased and all the memory cells are rewritten into identical data (\"0\" or \"1\"). Therefore, the conventional verify operation, in which memory cells of about 1 bit to 4 kbits are simultaneously verified by sense amplifiers similar to the write operation, is hard to be considered efficient.\nIn view of the above, a variety of erase verify methods as follows have been proposed. It is to be noted that the term \"verify\" mentioned hereinafter means the erase verify.\nThe discrimination between 0 and 1 of information held in the memory cells of the nonvolatile semiconductor storage device is normally executed by detecting whether a threshold voltage of the memory cell transistor is high or low. However, the verify methods described below are the verify methods of a memory cell array in which the low threshold voltage state is defined as an erased state. The verify operation in this case is required to collectively check the event that the threshold voltages of all the memory cell transistors have been reduced (i.e., the event that a current flows through all the memory cell transistors when a voltage intermediate between the high and low threshold voltages is applied to the gates of all the memory cell transistors).\n(1) The erase operation is completed at the point of time when the event that a current has flowed through the memory cell transistor exhibiting the minimum threshold voltage among all the memory cell transistors is detected (Japanese Patent Laid-Open Publication No. HEI 4-3395.)\n(2) The n(n: integer) memory cells connected to one word line are simultaneously verified by the same number of decision circuits (Japanese Patent Laid-Open Publication No. HEI 8-227590). This verify method has the same concept as that of the line test proposed in connection with the prior art DRAM (Dynamic Random Access Memory) and the like.\n(3) In a virtual ground type memory cell array, when a current flows by applying a voltage in series across the sources and the drains of a number of memory cell transistors connected to one word line, it is decided that all the memory cells connected to the word line have been erased (Japanese Patent Laid-Open Publication No. HEI 7-111901).\nHowever, the aforementioned prior art nonvolatile semiconductor storage device verify methods have the problems as follows.\nThat is, as stated before, the above verify methods are executed by collectively checking whether or not a current flows through all the memory cell transistors when the voltage intermediate between the high and low threshold voltages is applied to the gates of all the memory cell transistors.\nHowever, the operation of collectively verifying the event that a current flows through all the memory cells is required to detect one memory cell through which no current flows among a number of memory cells through which currents are flowing, and this is very hard to be achieved physically.\nFor example, according to the verify method listed in the item (1), the erase operation is executed by reducing the threshold voltages of the memory cell transistors from a programmed state D1 to an erased state D2 as indicated by the memory cell transistor threshold voltage distribution shown in FIG. 13. Then, the completion of the verify operation is determined by detecting the event that the threshold voltage A of the memory cell which is most likely to be erased has lowered from a word line selection voltage C in the verify operation. Accordingly, there still remains the possibility that the threshold voltage B of the memory cell which is least likely to be erased is greater than the selection voltage C and is kept intact in the programmed state, and this poses the problem that the erasing of all the memory cells has been verified.\nAccording to the verify method listed in the item (2), the n memory cells connected to one word line are simultaneously verified by the same number of decision circuits, and therefore, the problem of the verify method (1) can be solved. However, the method has the problem of an increase in area for the requirement of n decision circuits identical in number to the n memory cells to be verified collectively and the problem of an increase in time of the verify operation for the requirement of the verify operation times in a cycle identical to the number m(m: integer) of the word lines.\nAccording to the verify method listed in the item (3), the problem of the verify method (1) can be solved. However, if the ON-state resistance of the memory cell transistor and the substrate effect are taken into consideration, the method has the problem that the current is too infinitesimal and the problem that not so much memory cell transistors can be connected and simultaneously verified due to an increase in threshold value."}
-{"text": "(1) Field of the Invention\nThe present invention relates to a structure of a water-cooled engine, in particular relates to a snowmobile engine cooling arrangement having cooling water passages for cooling a snowmobile engine.\n(2) Description of the Prior Art\nConventionally, most small snow vehicle such as snowmobiles and the like use two-cycle engines because they are relatively simple in structure, light and compact and yet powerful. Recently, however, because of regulation of exhaust gas due to environmental issues and aiming at improvement in reduction of fuel consumption rate, there is a trend toward employing four-cycle engines.\nAs the engine structure for snowmobiles, water-cooled engines have been generally employed because they can stably cool the engine by cooling water, preventing the occurrence of overheating or overcooling and improving the output power and yet are effective in reducing noise.\nAs an example of conventional engine cooling arrangements, Japanese Patent Application Laid-Open No. 2001-12243 discloses a cooling system for a snowmobile engine. This cooling system is comprised of a cooling water jacket formed inside the engine, a heat exchanger for cooling the cooling water and a cooling water pump for transferring the cooling water so as to circulate the cooling water inside the engine, and further includes: a cooling water bypass through which the cooling water can pass, avoiding its flow into the heat exchanger; and a switching valve for permitting the cooling water to flow into the cooling water bypass, wherein the switching valve is adapted to switch the cooling water circuit in such a manner that the cooling water is allowed to pass through the heat exchanger when the cooling water is equal to or higher in temperature than a predetermined level, and the cooling water is restricted from passing through the heat exchanger and is conducted to the cooling water bypass when the temperature of the cooling water is lower than the predetermined temperature level.\nMeanwhile, in contrast to two-cycle engines, which are high in power, light and compact, four-cycle engines need a camshaft and oil lubrication, inevitably tending towards large size.\nIn particular, when a four-cycle engine is included in a snowmobile, it is necessary to provide a contrived layout of the oil pan configuration, intake and exhaust systems and associated auxiliaries, in order to make the body and engine hood have similar size to those of a two-cycle engine.\nFurther, in general, when the mounted engine is of a water cooled type, a reservoir tank should be interposed within the piping of the cooling water. That is, it is necessary to take into account the installation space for this reservoir tank. Further, the reservoir tank needs to be arranged at a site where it is easy for cooling water to be re-supplied.\nSince the snowmobile is used in cold areas, in order to stabilize the idling engine speed at the start of operation or in order to regulate the surging of the thermostat, it is necessary to consider the layout of the cooling water piping and the position of the thermostat.\nFurther, the cooling water pump for transferring cooling water is driven by a belt which is driven by the rotation of the crankshaft. This means the enlargement of the longitudinal size of the engine. Therefore it is necessary to arrange the cooling water pump at an appropriate position. Moreover, it is necessary to choose the position of the cooling water pump optimally, considering the positional relationship with an alternator, which is belt driven, so that the front-to-rear size of the engine will be minimized.\nThe present invention has been devised in view of the above conventional problems, it is therefore an object of the present invention to provide a cooling arrangement for a snowmobile engine, which is able to stabilize the engine idling speed at the start of operation and is improved in the work performance in supplying cooling water and realizes a space-saving engine layout.\nIn order to achieve the above object, the present invention is configured as follows:\nIn accordance with the first aspect of the present invention, a cooling arrangement for a snowmobile engine, includes: a cooling water jacket formed inside an engine; a heat exchanger for cooling water; a water pump for ejecting cooling water; and a thermostat for controlling supply of the cooling water ejected from the water pump into the engine, so as to cool the engine by circulating the cooling water inside the engine, and is characterized in that:\nthe engine is mounted with its cylinder head at top, and an exhaust manifold is disposed on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder head while an intake manifold is disposed on the rear side, with respect to the vehicle\"\"s direction of travel, of the cylinder head;\nthe water pump is disposed on the front side, with respect to the vehicle\"\"s direction of travel, of cylinder block of the engine, under the exhaust manifold and is coupled to a crankshaft projected from one side wall, with respect to the widthwise direction of the body, of the cylinder block, so that it is driven by the rotational force which is transmitted through a drive belt from a rotational member fitted on one end of the crankshaft;\na cooling water inlet port for leading cooling water into the cooling water jacket inside the engine is provided on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder block and under the exhaust manifold at a position close to the drive belt;\na cooling water outlet port through which cooling water is taken out from the interior of the engine is arranged on the engine\"\"s side face at a position over and opposite to the drive belt; and\nthe thermostat is arranged at the cooling water outlet port, at a point downstream with respect to the flow of cooling water.\nIn accordance with the second aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the cooling water outlet port is comprised of a thermo housing for incorporating the thermostat and a thermo cap for covering the thermostat, the thermo cap being connected to the heat exchanger by way of a cooling water passage.\nIn accordance with the third aspect of the present invention, the cooling arrangement for a snowmobile engine having the above second feature is characterized in that the thermo housing is formed with a first cooling water bypass passage which branches off at a point upstream, with respect to the flow of cooling water, of the thermostat and is connected by way of a throttle body to the cooling water passage at a point upstream of the water pump.\nIn accordance with the fourth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above second feature is characterized in that the thermo housing is formed with a second cooling water bypass passage which branches off at a point upstream of the thermostat and is connected directly to the cooling water passage at a point upstream of the water pump without passing through the throttle body.\nIn accordance with the fifth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above third feature is characterized in that the thermostat is used to control the outlet of cooling water after passage of the cooling water jacket inside the engine while the thermo housing side and cooling water passage side are kept so as to be in constant communication to each other through the first cooling water bypass passage and the second cooling water bypass passage.\nIn accordance with the sixth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above fourth feature is characterized in that the thermostat is used to control the outlet of cooling water after passage of the cooling water jacket inside the engine while the thermo housing side and cooling water passage side are kept so as to be in constant communication to each other through the first cooling water bypass passage and the second cooling water bypass passage.\nIn accordance with the seventh aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that, in the snowmobile, an oil filter is arranged on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder block and an oil cooling means is interposed between the oil filter and the cylinder block.\nIn accordance with the eighth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that, in the snowmobile, an oil cooling means is provided on the front side, with respect to the vehicle\"\"s direction of travel, of a crawler for causing the snowmobile to move.\nIn accordance with the ninth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that, in the snowmobile, a muffler is disposed in front of the engine body in the engine room and an oil cooling means is arranged in front of the muffler.\nIn accordance with the tenth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the water pump is arranged between the engine and the exhaust passage provided in front of the engine.\nIn accordance with the eleventh aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the water pump is arranged over an oil pan.\nIn accordance with the twelfth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the water pump is arranged under the exhaust manifold.\nIn accordance with the thirteenth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the engine has an alternator provided at a position opposite the water pump with the cylinder block in between, the three components being laid from the front to rear with respect to the vehicle\"\"s direction of travel.\nIn accordance with the fourteenth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that a cooling water reservoir tank is arranged in the rear of the engine, at the side of the intake manifold, over an oil pan, heat exchanger and alternator.\nIn accordance with the fifteenth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that a cooling water reservoir tank is arranged without being projected outwards beyond a drive belt for driving an alternator and water pump, when viewed from top.\nAccording to the present invention, the following effects can be obtained.\nFirst, in a cooling arrangement for a snowmobile engine, the engine is mounted with its cylinder head at top, and an exhaust manifold is disposed on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder head while an intake manifold is disposed on the rear side, with respect to the vehicle\"\"s direction of travel, of the cylinder head. The water pump is disposed on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder head of the engine, under the exhaust manifold and is coupled to a crankshaft projected from one side wall, with respect to the widthwise direction of the body, of the cylinder block, so that it is driven by the rotational force which is transmitted through a drive belt from a rotational member fitted on one end of the crankshaft. Further, a cooling water intake port for leading cooling water into the cooling water jacket inside the engine is provided on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder block and under the exhaust manifold, at a position close to the drive belt while a cooling water outlet port through which cooling water is taken out from the interior of the engine is arranged on the engine\"\"s side face opposite to the drive belt side or on the opposite side at a position above the crankshaft from which engine power is taken. This arrangement makes it possible for the cooling water to flow through the engine interior, approximately diagonally from the bottom to the top, hence the water is able to efficiently flow and be discharged without stagnation inside the cylinder head and inside the cylinder block. Accordingly, this configuration provides an engine of an improved cooling efficiency compared to the conventional configuration. Further, since supply and discharge of cooling water is achieved utilizing the dead space near the water pump and under the exhaust manifold, this arrangement realizes an engine of a highly efficient space usage.\nSince the cooling water outlet port is formed with a thermo housing incorporating a thermostat and a thermo cap covering the thermostat while the thermo cap is connected to the heat exchanger by way of a cooling water passage, e.g., a cooling water hose so as to control the cooling water at the cooling water outlet, by making a so-called outlet control, it is possible to perform correct cooling water control at a stable cooling water temperature.\nSince the thermo housing is formed with a first cooling water bypass passage which branches off at a point upstream, with respect to the flow of cooling water, of the thermostat and is connected by way of a throttle body to the cooling water passage at a point upstream of the water pump, it is possible to prevent icing and make stable intake control by maintaining the cooling water supplied to the throttle valve at an approximately fixed temperature.\nSince the thermo housing is formed with a second cooling water bypass passage which branches off at a point upstream, with respect to the flow of cooling water, of the thermostat and is connected directly to the cooling water passage at a point upstream, with respect to the flow of cooling water, of the water pump without passing through the throttle body, the cooling water warmed by the engine is supplied to the engine. This makes it possible to regulate the temperature of the cooling water at an approximately constant level under the thermostat control of supply of the cooling water cooled through the heat exchanger, hence realizing a good operational condition of the engine.\nSince control of cooling water by the thermostat is performed by regulating the outlet of cooling water after passage of the cooling water jacket inside the engine, it is possible to make reliable cooling water control based on the cooling water of a stabilized water temperature. Further, since the thermo housing side and cooling water passage side are kept so as to be in constant communication to each other through the first cooling water bypass passage and the second cooling water bypass passage, this makes it possible to regulate the temperature of the cooling water in the engine at an approximately constant level. As a result, the thermostat for opening and closing the main cooling water passage in response to the cooling water temperature quickly reacts to the cooling water temperature that varies dependent on the operational state of the engine, to thereby prevent the occurrence of engine seizure and other troubles.\nSince, in the snowmobile, an oil filter is arranged on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder block and an oil cooling means, e.g., an oil cooler, is interposed between the oil filter and the cylinder block, it is possible to efficiently cool not only the engine but also the engine oil, thus achieving a beneficial operational condition of the engine.\nSince, in the snowmobile, an oil cooling means is provided on the front side, with respect to the vehicle\"\"s direction of travel, of a crawler for causing the snowmobile to move, it is possible to perform improved cooling by cooling oil with the scattered snow. Further, since two divided heat exchanging means or heat exchangers are provided, where one is used as usual for cooling water and the other for cooling oil, it is possible to realize different cooling functions with a simple structure.\nSince a muffler is disposed in front of the engine body in the engine room in the snowmobile and an oil cooling means is arranged in front of the muffler, this arrangement allows the oil cooling means to receive the flow of air, during travel, ahead of the muffler and exhaust pipe which are higher in temperature than the oil cooler. Thus, it is possible to cool the oil at an improved efficiency.\nSince the water pump is arranged between the engine and the exhaust passage provided in front of the engine, it is possible to realize a space-saving engine layout utilizing the dead space under the exhaust passage.\nSince the water pump is arranged over an oil pan, it is possible to prevent the engine height from increasing.\nSince the water pump is arranged under the exhaust manifold, it is possible to realize a space-saving engine layout utilizing the dead space under the exhaust manifold.\nSince the engine has an alternator provided at a position opposite the water pump with the cylinder block in between, the three components being laid from the front to rear with respect to the vehicle\"\"s direction of travel, it is possible to realize a space-saving and well-balanced drive belt layout.\nSince a cooling water reservoir tank is arranged in the rear of the engine, at the side of the intake manifold, over an oil pan, heat exchanger and alternator, this arrangement of the reservoir tank at the top of the engine room makes it easy for cooling water to be re-supplied and other maintenance.\nSince a cooling water reservoir tank is arranged without being projected outwards beyond a drive belt for driving an alternator and water pump, when viewed from top, this arrangement realizes a compact engine configuration."}
-{"text": "Aluminium profiles are used extensively in applications such as window and door frames, display counters and partitions. Such aluminium profiles come in several finishes, among which are metallic finishes such as the well-known silver-metallic and gold or bronze metallic finishes. Recently a brown-metallic finish has become popular as well. These finishes involve anodizing the aluminium. Alternatively, aluminum profiles have been coated with paints which are baked onto the aluminum. To anodize aluminum, it has to undergo an electro-chemical process which is rather complex, expensive and limits the colours obtainable. Aluminum profiles that are painted require first to undergo a preliminary process of undercoating in order to provide a suitable surface for the paint to adhere to, since paints do not adhere well to aluminum.\nAlthough many advances and developments have taken place in the coating industry, including the art of finishing aluminum profiles, there are a number of problems that have as yet not been overcome. First of all the finished aluminum profiles, which are used for so many diverse applications, require special resistance against degradation by light, chemicals, corrosive atmospheric gases and other pollutants, because they are used so much under outdoor conditions in displays and window frames where ultra-violet light, industrial gases and other weather conditions deteriorate the coating and the aluminum. This is particularly true for use of aluminum near coastal areas, where oxygen and salt water in the air subjects the profiles to severe corrosion.\nA further problem is caused by the fact that modern household and sanitary cleaning materials contain bleaches and alkali which react with aluminum and deteriorate its surface. These problems dictate that the aluminum profiles for use in such applications must have complete protection from contact between any corrosive atmosphere and the aluminum. This protection is very difficult to achieve in profiles with intricate channels and grooves, as is the case with aluminum profiles. Yet the demand for aluminum profiles for such applications is ever-increasing, particularly for profiles with metallic finishes at a reasonable price."}
-{"text": "1. Field of the Invention\nThe invention relates to a connector.\n2. Description of the Related Art\nU.S. Pat. No. 7,448,908 discloses a waterproof connector with first and second housings that are fit together by operating a lever. The first housing has a frame with concave accommodation parts and sub-connectors having terminal fittings are accommodated in the concave accommodation parts. A collective rubber stopper is accommodated in the concave accommodation parts to seal gaps between electric wires pulled out of the sub-connectors and inner peripheries of the concave accommodation parts. An engaging part is formed on the outer surface of the frame and defines a supporting shaft for pivotably supporting the lever. The lever has a cam groove that engages a cam follower of the second housing. Pivoting the lever generates a cam action between the cam groove and the cam follower that urges the first and second housings together.\nFit-on resistance between the first and second housings creates a reaction force that acts on the engaging part when a pivotal operation force is imparted to the lever. Thus, there is a fear that the outer wall of the frame will deform. Such a deformation will generate a gap between inner surfaces of the concave accommodation portions and the outer surface of the collective rubber stopper and could adversely affect the waterproof performance.\nThe invention has been completed based on the above-described situation. It is an object of the invention to prevent deterioration of the waterproof performance of a connector."}
-{"text": "The present invention relates to a refrigerator unit and/or a freezer unit having one or more storage shelves for the reception of refrigerated goods or frozen goods, with the appliance having one or more extension arrangements which are connected to the storage shelves and by means of which the storage shelves can be pulled out of and pushed into the appliance.\nPull-out storage shelves of refrigerator units and/or freezer units have the substantial advantage with respect to fixedly mounted, non-movable shelves that the refrigerated goods or frozen goods in the rear region of the shelves are comparatively easily accessible even when refrigerated goods and/or frozen goods are located in the front region of the shelves. A clearing away of the front region is thus not absolutely necessary. A disadvantage of previously known storage shelves consists of the fact, however, that the position of the shelves is fixedly set so that the appliances are comparatively inflexible with respect to the adaptability of the inner space to changing needs."}
-{"text": "Semiconductor devices are subjected to a series of test procedures in order to confirm functionality and yield, and to assure quality and reliability. This testing procedure conventionally includes \"probe testing\", in which individual dice, while still on a wafer, are initially tested to determine functionality and speed. Probe cards are used to electrically test die at that level. The electrical connection interfaces with only a single die at a time in wafer, before the dice are singulated from wafer.\nIf the wafer has a yield of functional dice which indicates that quality of the functional dice is likely to be good, each individual die is traditionally assembled in a package to form a semiconductor device. Conventionally, the packaging includes a lead frame and a plastic or ceramic housing.\nThe packaged devices are then subjected to another series of tests, which include burn-in and discrete testing. Discrete testing permits the devices to be tested for speed and for errors which may occur after assembly and after burn-in. Burn-in accelerates failure mechanisms by electrically exercising the devices (devices under test or DUT) at elevated temperatures, elevated dynamic biasing schemes. This induces infant mortality failure mechanisms and elicit potential failures which would not otherwise be apparent at nominal test conditions.\nVariations on these procedures permit devices assembled onto circuit arrangements, such as memory boards, to be burned-in, along with the memory board in order to assure reliability of the circuit board and the circuit board assembly and manufacturing process, as populated with devices. This closed assembly testing assumes that the devices are discretely packaged in order that it can then be performed more readily.\nSemiconductor packaging has been referred to in terms of \"levels\" of packaging. The chip capsule generally constitutes a first level of packaging. A second level would then be a \"card\" or a printed circuit board. A third level may include second level packaging combined with a motherboard. A fourth level may follow the third level. In each case, the packaging to any level involves cost.\nIt is proposed that devices be packaged without conventional lead frames. This creates two problems for being conventional test methods. Firstly, discrete testing is more difficult because the conventional lead frame package is not used. Furthermore, multiple device may be assembled into a single package, thereby reducing the performance of the package to that of the die with the lowest performance. This is because the ability to presort the individual dice is limited that obtained through probe testing. Secondly, the packaging may have other limitations package assembly defect mechanisms which are aggravated by burn-in stress conditions so that the packaging becomes a limitation for burn-in testing.\nAccording to the invention represented by U.S. Pat. No. 4,899,107, to Alan Wood and Tim Corbett, a reusable burn-in/test fixture for discrete die is provided. The fixture consists of two halves, one of which is a die cavity plate for receiving semiconductor dice as the devices under test (DUT); and the other half establishes electrical contact with the dice and with a burn-in oven.\nThe first half of the test fixture contains cavities in which die are inserted circuit side up. The die will rest on a floating platform. A support mechanism under the die platform will provide a constant uniform pressure or force to maintain adequate electrical contact to the die contacts on the DUT to probe tips on the second half. The support mechanism will compensate for variations of overall die thickness.\nThe second half has a rigid high temperature rated substrate, on which are mounted probes for each corresponding die pad. Each probe is connected to an electrical trace on the substrate (similar to a P.C. board) so that each die pad of each die is electrically isolated from one another for high speed functional testing purposes. The probe tips are planar so that contact to each die pad occurs simultaneously. The probe tips are arranged in an array to accommodate eight or more dice. The traces from the probes terminate in edge fingers to accept a conventional card edge connector. The geometry of the probes and edge fingers are optimized to avoid electrical test artifacts.\nThe two halves of the test fixture are joined so that each pad on each die aligns with a corresponding electrical contact. The test fixture is configured to house groups of 8 or 16 die for maximum through put efficiency of the functional testers. The test fixture need not be opened until the burn-in and electrical test are completed. After burn-in stress and electrical test, the die are removed from the test fixture and depositioned accordingly. The fully burned-in and tested die are available for any type of subsequent assembly applications.\nThat technique allows all elements of the burn-in/test fixture to 100% reusable, while permitting testing of individual dice in a manner similar to that accomplished with discrete packaged semiconductor devices.\nAn ability to extend accelerated burn-in and functional/parametric/speed testing of dice to include accelerated burn-in and functional, parametric and speed testing while the dice are still on the wafer would have several advantages. Since each step in the assembly and package process represents commitment of resources, early determination of defective parts or ability to predict a failure at a conventional burn-in stage is advantageous. It would be further advantageous to be able to predict a failure at a burn-in stage prior to assembly. Clearly, if a part can be made to fail prior to assembly, assembly resources can be directed to a higher percentage of good parts.\nThere exists a significant market for uncut fabricated wafers. These wafers are referred to as \"probe wafers\" because they are delivered after probe testing, which follows fabrication. The purchase of probe wafers is primarily by \"ASIC assembly houses\" which custom package parts, including parts traditionally considered to be \"commodity\" chips. The purchase of uncut wafers is usually based on the recent yield rate of the semiconductor manufacturer, but recent yields are not a strong indicator of the yield of any given wafer lot. Furthermore, the assembly process techniques used by the assembly house have a significant effect on yield.\nCharacterization, such as speed grading is even more variable than yield. While a packaged DRAM is purchased by the consumer based on the part's speed grade, speed grading of probe wafers is almost a matter of conjecture. That means that it is happenchance as to whether the assembly house purchases a wafer of mostly \"-10\" parts (100 ns) or mostly \"-6\" parts (60 ns).\nRecent developments in fabrication technology have resulted in such speed characterizations being more uniform on any given wafer. This has made it possible to provide wafers in which a majority of good dice have speed grades which do not greatly exceed an average for the wafer. Such uniformity, along with an ability to achieve fuse repairs and patches, have made wafer scale integration of arrays and cluster packaging practical.\nOther developments include an ability to track individual dice on wafers, starting at probe. Traditionally, probe identifies bad dice, (example, an ink spot.) The assembly process is continued only for dice which do not have the ink spots. By computer tracking, the ink spot becomes superfluous, as a map of good and bad dice are stored and transferred to subsequent assembly steps.\nAlthough the dice are singulated, there are cases in which the discrete parts are reassembled into an array after assembly. An example is in computer memory, in which one or more banks of memory are composed of arrays of memory chips. It would be advantageous to be able to select good dice on a wafer and assemble the dice into an array without singulating the dice. This would allow a much denser array of good clustered dice on a single piece of silicon.\nIt is an object of the invention to increase handling efficiency, while at the same time reducing the required size of the test fixture."}
-{"text": "The present day digital camera includes many functions. For example, there is known a digital camera which has a character photographing mode in addition to a normal photographing mode. The normal photographing mode is the mode in which ordinary snapshots, for example, photographs of scenery are taken. In the character photographing mode, the digital camera subjects data to image processing so as to improve the appearance of the characters.\nThere is a digital still camera disclosed in Japanese patent application laid open No. HEI 08-125870. This digital still camera has a structure as follows. The digital still camera has the normal photographing mode and the document photographing mode. Any one of these two modes can be selected. When the normal photographing mode is selected, image data is compressed by a natural picture compression unit. Whereas image data is compressed by a document compression unit (a compression method appropriate for an image with a small number of gray levels in achromatic color) when the document photographing mode is selected. The compressed image data is then stored in an externally provided memory. Therefore, the image data can efficiently be compressed irrespective of whether the data belongs to a picture or belongs to a photograph of a document without degradation in image quality.\nHowever, the conventional type of digital camera has a problem that the data of the captured image can not effectively be used. More specifically, the digital camera having the character photographing mode can not transmit the captured image as through a facsimile. In order to transmit the image as a facsimile, following process is carried out. First, the image is transmitted to a computer where the image is subjected to processing such as expansion, image conversion, compression, or the like and then only the image can be transmitted as a facsimile.\nFurther, in the digital still camera described above, the compression method is uniquely determined based on the photographing mode. For example, when a document (for example, an encyclopedia) which contains picture of scenery is photographed, and if the picture is photographed in the document mode, the image data is compressed by the compression method suitable for an image with a small number of gray levels in achromatic color. Resultantly, the picture of the scenery can not be reproduced properly. Therefore, the user is required to photograph the same document in the normal photographing mode. Thus, the user has to choose the photographing modes based on the purpose of using the document."}
-{"text": "This invention relates to fiber optic devices and, more particularly, to devices for monitoring, switching, attenuating or distributing optical signals in fibers.\nOptical transmission systems are being developed to handle large amounts of communications traffic. One type of optical transmission system uses optical fibers as the transmitting media. Optical fibers with sufficiently low loss for practical long distance optical communication systems have been developed only recently.\nAny communication system requires both transmitting and receiving apparatus. It is desirable that the transmitting source have a constant power output that is independent of temperature changes and component aging. One source of optical power is a laser which changes output power in response to changes of both temperature and aging. Output power of the laser can be stabilized by monitoring the light actually coupled into the fiber and by using such light to develop a signal for feedback control of the laser. To effectively monitor such light in the fiber for controlling the laser, a tap should monitor the transmitted light without altering the characteristic of the transmitted power. All transmitted modes should be monitored equally because any substantial preference of one mode or modes over another mode or modes is unsuitable for accurate feedback control of the laser. Available fiber optic taps substantially prefer one optical mode or modes over another mode or modes.\nFiber optic systems need optical tap devices for economically monitoring and controlling the optical power. A wide variety of devices are being developed to fill these needs. Development of a device may result in a specific and perhaps optimum component for each particular application. Such a specific development can be costly. As an alternative a designer may desire to convert the optical signal to an electrical signal and process the latter electronically taking advantage of well developed and versatile electronic technology. This electronic approach introduces noise into the system as well as requiring a source of external power.\nProblems created by developing specific components for each application and by converting optical signals to electronic signals for processing are avoided by the use of an optical fiber access port in a tap for monitoring and controlling optical power and for building many optical control devices without substantial further development effort regardless of changes introduced by new fiber materials and structures.\nA practical optical communications system requires at least a few different optical devices for establishing desirable interconnection arrangements. Optical switches can be used so that part of the system can be rearranged conveniently either for testing or for providing a variety of communications needs. Additionally, a variable optical attenuator can be used for adjusting the magnitude of power being transmitted in any part of the system. Such attenuators are helpful in setting levels for testing and for establishing some branching arrangements. Another device which facilitates the design of optical communication systems is a multibranch distributor. Such a distributor divides incident light, transmitted along one fiber, into portions for further transmission along two or more other fibers.\nAs in the case of the optical fiber tap, the variable attenuator and the multiple branch distributor should perform their designated functions without substantial preference of some mode or modes over another mode or modes. Accurate transmission of information often requires faithful communication of the input signal characteristic throughout the system.\nAvailable fiber optic devices are relatively complicated and expensive devices whereas it is desirable for reasons of cost and ease of operation to have simple, inexpensive fiber optic devices. Such simple, inexpensive devices enhance the practicality of optical communications systems arrangements.\nThese problems of complicated design, high cost and mode selectivity of the existing fiber optical devices are solved by a group of fiber optic devices, each including an access port.\nIt is an object to provide a family of simple fiber optic devices for coupling optical power between fibers.\nIt is a further object to provide fiber optic devices which are substantially insensitive to optical modes carried by a fiber.\nIt is another object to provide a basic unit for constructing different optical control devices."}
-{"text": "The present invention relates to a microorganism capable of producing L-lysine and a method for producing L-lysine thereby. More specifically, the present invention relates to a strain possessing resistance to arginine analogues and other analogues besides the characteristics known as necessary for the production of L-lysine, and a method for producing L-lysine comprising culturing the strain in a nutrient culture medium and recovering L-lysine from the culture broth.\nL-Lysine, which is one of essential amino acids, has been used as an animal feed supplement because of its ability to improve the quality of feed by increasing the absorption of other amino acids, and used as a food supplement as well. L-Lysine has been used also in medicine, particularly as ingredients of infusion solutions.\nL-Lysine can be produced by fermentation using an auxotrophic mutant and/or regulatory mutant, or by enzymatic conversion which comprises treating DL-.alpha.-amino-.beta.-caprolactam (hereinafter referred to as ACL) with L-ACL-hydrolase and ACL-racemase. For commercial production of L-lysine, however, fermentation processes have been more predominantly employed. For example, Japanese Patent Laid-Open No. 82-9797, and Japanese Patent Publication Nos. 82-14839 and 86-35840 disclose the methods using strains having characteristics of vitamin B.sub.1, pantothenic acid and/or biotin requirement for growth; resistance to lysine, threonine and/or isoleucine analogues; resistance to one or more antibiotics e.g. bacitracin, penicillin G and/or polymyxin."}
-{"text": "1. Field of the Invention\nThe present invention relates to a method for assembling a camera module for a vehicle.\n2. Description of Related Art\nCamera modules are being increasingly used these days in vehicles, for example for use in night vision systems or as lane departure warning (LDW) systems. Particular demands in terms of robustness are made on such cameras with regard to their installation in the vehicle, given the vibrations that unavoidably occur during operation.\nTwo fundamentally different approaches are known for the assembly of camera modules for vehicles: An imager module is used which encompasses an image sensor chip, a basic housing, and a lens, the lens being produced and tested as an independent module. Imaging modules of this kind can be attached to the circuit board, for example by thermal joining, e.g., using soldering methods, or by way of a cable end having a plug. The circuit board, completed in this fashion, is then installed into the housing, which must have a corresponding receptacle for the lens. A disadvantage of this procedure is the fact that imager modules having chip and lens dimensions suitable for motor vehicle applications cannot be reliably produced using common soldering methods (e.g., reflow methods). Because of the high heat capacity, poor solder connection points can form, or the optics can be damaged. On the other hand, the plug contacts known as alternative solutions are not fault-free over the long term in vehicles. A further problem that has emerged here is sealing of the entire camera module in the context of a circuit board plus imager module configuration, since mechanical stresses acting on the imager module often occur; this can result in subsequent damage.\nA further assembly concept is based on configuration of the circuit board with an independently packaged image sensor; this typically involves a ceramic housing having a glass cover. The sensitive pixel array can thereby be sufficiently protected from particles. Such housings can moreover be processed using reflow processes. The lens can be focused either in a lens holder mounted on the circuit board, or in the housing shell of the camera module. This approach, however, exhibits a very long tolerance chain because of the larger number of individual components requiring installation; after assembly this can result, in the least favorable case, in inaccurately operating camera modules due to addition of the greatest tolerances.\nPublished German patent document DE 10 2004 001 698 discloses an optical module. According to this approach, an optical module encompasses an image acquisition device made up of optical and electronic components, the components of the image acquisition device being disposed together on a carrier plate. The components of the image acquisition device are disposed on a first principal surface of the carrier plate, an optical system associated with the image acquisition device being disposed on a second principal surface of the carrier plate. The carrier plate and the optical system are embedded in a sealing compound. The optical module encompasses a holding element for securing the optical module to a vehicle part or the like. The holding element is sealed to the optical module, and the optical module encompasses a shield against electromagnetic interference radiation. According to this approach, the shield can also be embedded in the sealing compound. The shield is preferably embodied as a net or lattice embedded in the sealing compound, or can be made up of particles embedded in the sealing compound."}
-{"text": "The invention relates to a method and a plant for winding a ship into a coil and applies in particular to the winding of a steel band after hot rolling.\nWhen exiting out of a hot band train, the steel band is carried generally on a roller table, towards one or several coiling machines situated at a certain distance downstream of the roller table, generally under it.\nUpstream of the coiling machines, the band is driven by a clamping device, followed by guiding means which direct the band towards a rotary chuck whereon the band, driven by friction, is wound in superimposed turns.\nFor winding a new band, the head of the band is applied on the external face of the rotary chuck in order to be driven by friction downstream. Then, the band is guided by a curved bending plate so as to be bent while remaining applied on the chuck, along an angular sector which increases gradually, which enhances the friction driving effect. After one revolution, the head of the band passes below the internal face of the hand to form a second turn and so on. In order to keep the tension necessary for winding, the band is held upstream by a device including, usually, a deflecting roll and a pinch roll driven into rotation at an angular velocity slightly smaller than that of the chuck, wherein the band may thus be stretched along a theoretical winding plane tangent, upstream, to the deflecting roll, along a generatrix thereof, and, downstream, to the coil being winding. Usually, the chuck is retractable so as to reduce in diameter to enable the removal of the coil, once it has reached the requested diameter.\nThe document U.S. Pat. No. 2,918,226, for instance, describes a winding plant of such type including, downstream of a deflecting roll placed at the outlet of a roller table, a rotary chuck arranged inside a pit and associated with two winding assemblies, each comprising an application roll of the band on the cylindrical external face of the chuck and a curved bending plate for guiding the band during the winding process. In the arrangement of document U.S. Pat. No. 2,918,226, the assembly is placed on a supporting chassis mounted so that it is capable of being pivoted around an axis parallel to the axis of the chuck and actuated by a jack, so as to move between a spaced apart position and a working position for which the bending plate is substantially parallel to the external face of the chuck. On the other hand, the application roll is mounted itself on an auxiliary chassis hinged on the supporting chassis of the bending plate and pushed towards the chuck by a spring so as to exert an elastic force to apply the band on the chuck, along a pushing direction running through a back-up generatrix of the external face of the chuck and the axis thereof.\nMoreover, the application roll can be driven into rotation around its axis in a reverse direction relative to the direction of rotation of the chuck, in order to apply, to the external face of the band, an additional driving friction force which is added to the main driving friction force exerted by the chuck on the internal face of the band.\nAt the beginning of the winding process, the head of the band may rest on the bending plate situated downstream of the application roll and is guided in the annular gap between the bending plate and the external face of the chuck so as to be winded around the chuck. After one revolution, the head of the band may run below the hand so as to form a second turn, the application roll, pushed by a spring, being capable to move apart slightly.\nIn a more developed embodiment described in document U.S. Pat. No. 3,587,274, an additional roll is mounted, upstream of the application roll and of the chuck, on a chassis hinged around an axis, in order to exert a pre-bending torque on the band.\nUntil now, such devices were satisfactory since they were used for winding steel bands with a relative small thickness, for instance ranging between 2 and 12 mm.\nFor such thicknesses, indeed, the band bends easily and may be guided in the annular gap provided between the chuck and the bending plate, so as to be applied on the cylindrical face of the chuck in order to be driven by friction.\nHowever, it is now contemplated to use devices of such type for winding hot bands with increased thicknesses, up to for example 25 mm. Moreover, one may have to wind products with high elastic limits, for instance of the order of 370 MPa at the temperature of use.\nStill, it has appeared that, for relatively high thicknesses, the usual arrangement of the type represented on FIG. 1 does not enable sufficient bending of the band, in particular at the beginning of its winding process.\nAt that time, indeed, the band is applied only on a small angular sector of the chuck, and the main driving force exerted by the chuck on the internal face of the band depends essentially of the pressure exerted by the application roll of the band to the chuck.\nHowever, in the previous arrangements disclosed by the documents mentioned above, this pushing force, exerted by a spring system or an elastic stop, would be insufficient to enable the winding of a hot band which is relatively thick and rigid."}
-{"text": "This section introduces aspects that may help facilitate a better understanding of the disclosure. Accordingly, these statements are to be read in this light and are not to be understood as admissions about what is or is not prior art.\nChronic pain is a major health concern that costs the US more than $635 billion per year (Gaskin and Richard, (2012) J. Pain 13:715-724). In addition to the financial impact, patients with chronic pain suffer extreme physical, emotional, and social burdens. For example, individuals often become socially isolated and confined to home as a result of their chronic pain that is not well-controlled by today's available treatments. The drugs used for the management of chronic pain include opioid analgesics, neuronal stabilizers such as anticonvulsants, and antidepressants. Opioids are the most widely used, and a recent NIH report indicates that there are significant problems associated with long-term opioid therapy for chronic pain (Volkow and McLellan (2016) N Engl J Med 374:1253-1263.). None of the agents provide sufficient relief to allow patients to return to their normal activity level. Moreover, current pharmaceutical industry has retreated from studying novel pain therapeutics due to the enormous risk (Skolnick and Volkow (2016) Neuron 92:294-297.). These observations indicate an essential need to identify new agents acting on unique targets in the war on chronic pain.\nAdenylyl cyclases are enzymes that catalyze the production of cyclic adenosine monophosphate, an important biological messenger. There are nine different membrane-bound human adenylyl cyclase isoforms, each with slightly different attributes. In particular, adenylyl cyclase 1 (AC1) is highly expressed in the hippocampus as well as in regions of the brain associated with pain. AC1 knock out animals show reduced neuropathic and inflammatory pain. Additionally, AC1 knock out mice show less reward when given opioids and show reduced symptoms of opioid dependence during withdrawal. Additional reports suggest that AC1 inhibition may also provide a useful therapeutic intervention for alcohol use disorder and autism (Bosse K E et al., J. Pharmacol. Exp. Ther. 2017, 363 (2) 148-155; Sethna F., et al. Nat. Commun. 2017, 8, 14359).\nUnfortunately, until now, the selective inhibition of ACs has not been achieved, and simultaneous inhibition of multiple adenylyl cyclase isoforms would likely result in significant adverse effects. There are unmet needs for better and safer medications targeting adenylyl cyclases for various therapeutic uses."}
-{"text": "1. Field of the Invention\nThe invention relates to methods and systems for detecting access points, and particularly to methods and systems for detecting rogue access points and a device for identifying rogue access points.\n2. Related Art\nWhen an organization such as a company uses a traditional wired network communication and computer system, it is generally not easy for network intruders to gain access into the company's network. For instance, the network intruders may need to gain physical access inside the company's building, and then find an Ethernet port from where they can electronically obtain confidential information of the company. This physical boundary provides a certain level of protection against network intruders.\nWhen a company uses a wireless network, the network intruders do not necessarily have to physically get inside the company's building in order to access the company's network. The network intruders only need to be present within the scope of the wireless network, and then gain access to the company's confidential information through a wireless connection.\nOne form of wireless security breach occurs when network intruders put an access point (AP) near to or in a company's building. An employee of the company may unknowingly connect the unauthorized access point to the company's network. The network intruders can then try to gain access to the employee's confidential information via the unauthorized access point. This kind of unauthorized access point is commonly known as a \u201crogue access point\u201d or \u201crogue AP.\u201d Therefore, a method for detecting rogue access points is needed in order to protect an organization's confidential information."}
-{"text": "1. Field of the Invention\nThe present invention relates to a plasma reactor and process for the production of hydrogen-rich gas from light hydrocarbons.\n2. Description of the Related Technology\nImproving the efficiency of energy production remains an important technological goal, owing to the significant economic benefits that result in almost every sector of the economy. One potential method for improving the efficiency of energy production is to provide an energy efficient method of converting light hydrocarbons to hydrogen-rich gas, to thereby increase energy production from natural gas.\nPlasma fuel converters such as plasmatrons are known to reform hydrocarbons to produce hydrogen-rich gas. DC arc plasmatrons, for example, are disclosed in U.S. Pat. Nos. 5,425,332 and 5,437,250. DC arc plasmatrons generally operate at low voltage and high current. As a result, these plasmatrons are particularly susceptible to electrode erosion and/or melting. DC arc plasmatrons also require relatively high power inputs of 1 kW or more and relatively high flow rates of coolant to keep the temperature in check.\nOther conventional methods for the conversion of light hydrocarbons to hydrogen-rich gas are generally energy inefficient and, as a result, in many small-scale applications, such as the production of hydrogen for fuel cells, the cost of hydrogen gas made by these methods is not competitive. Thus, there is a need in the art for a more energy efficient process for the conversion of light hydrocarbons to hydrogen-rich gas.\nU.S. Pat. Nos. 5,993,761 and 6,007,742 (Czernichowski et al.) describe processes for the conversion of light hydrocarbons to hydrogen-rich gas using gliding arc electric discharges in the presence of oxygen and, optionally, water. In the process, two electrodes having flat sheet geometry are arranged for arc ignition and subsequent gliding of the arc. The distance between the cathode and anode gradually increases to a point that no longer supports the gliding arc. As a result, the gliding arc disappears at one end of the electrodes, creating pulsed plasma wherein the properties of the plasma change with time. Due to the use of pulsed plasma, the process is relatively unstable over time. Reagents and oxygen are preheated using an external heat source. As a result of the preheating of the reagents and oxygen using an external heat source, the process suffers from poor energy efficiency. A premixed feed gas including hydrocarbons and oxygen is introduced to the reactor located at the central axis of the reactor.\nU.S. Pat. No. 5,887,554 (Cohn et al.) also discloses a system for the production of hydrogen-rich gas from light hydrocarbons. The system includes a plasma fuel converter for receiving hydrocarbon fuel and reforming it into hydrogen-rich gas. The plasma fuel converter can be operated using either pulsed or non-pulsed plasma and can utilize arc or high frequency discharges for plasma generation. Products from the plasma fuel converter are employed to preheat air input to the fuel converter. In one embodiment shown in FIG. 6, residence time in the reactor is increased by providing a centralized anode and a plurality of radial cathodes to thereby cause the arc to glide towards the center of the reactor under the influence of gas flowing in the same direction as the gliding arc.\nU.S. Pat. No. 6,322,757 (Cohn et al.) discloses a plasma fuel converter which employs a centralized electrode and a conductive reactor structure which acts as the second electrode for creation of a plasma discharge. Reagents are fed to the reactor just below the smallest gap between the electrodes and flow in the same direction as the gliding arc to thereby produce hydrogen-rich gas. In alternative embodiments, air and/or fuel are preheated by counter-flow heat exchange with the products of the reforming reaction and fed to the reactor either above or just below the smallest gap between the electrodes.\nAlthough some improvements in the energy efficiency of plasma fuel converters have been achieved, there remains a need for higher energy efficiencies for use of non-equilibrium low temperature plasma."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to golf equipment accessories and, more specifically, to a golf club protection device which can be inserted into the readily available individual golf club tubes which are sold and inserted into a golf bag to segregate and protect individual golf clubs. While these tubes do segregate the golf clubs one from the other they do not prevent the shafts from scraping the lip of the tube while being extracted from or inserted into during the course of play. This contact with the lip of the tube will remove the coating on graphite shafts.\nThe present invention comprising a flanged insert seats into the top opening of existing tubes by means of compression fitting or is therein permanently affixed by adhesive means and said flanged insert has an acrylic soft pile material circumferentially affixed to the interior wall of said flanged insert.\nAnother embodiment is provided having all of the properties of the preferred embodiment in addition to a club head bearing member comprised of a Y-shaped member fixedly attached to the exterior surface of the flanged insert therein providing means for supporting the club head within the V-portion and preventing any lateral movement of the club head which normally occurs during the course of play.\n2. Description of the Prior Art\nThere are other golf club shaft protection devices. Typical of these is U.S. Pat. No. 4,938,349 issued to Burns on Jul. 3, 1990.\nAnother patent was issued to Henry et al. on Jan. 4, 1994 as U.S. Pat. No. 5,275,278. Yet another U.S. Pat. No. 5,632,690 was issued to McConville on May 27, 1997 and still yet another was issued on Feb. 24, 1998 to King et al. as U.S. Pat. No. 5,720,388.\nA golf club protector for protecting the shaft of a club including a golf club housing tube having a protective interior and a protective collar provided at the upper end of the tube. Alternatively, an optional protective sleeve insertable within the tube having a protective collar attached thereto may be removably provided as a unit for a conventional tube. Where the housing tube includes a fixed collar and protective interior, the length of the tube may be predetermined in accordance with the length of a particular club shaft, or the tube may be provided with spaced apart cut markings at its lower end to facilitate shortening of the tube by the consumer.\nA golf club shaft protector is disclosed as including a hollow plastic tube of predetermined diameter and length with open upper and lower ends. The length of the hollow plastic tube substantially corresponds to the length of a golf club shaft and can be pre-selected and cut to the appropriate length. A flexible restricted throat element extends across the open upper end of the hollow plastic tube for resilient deformation upon the insertion of a golf club handle to allow passage of the golf club handle and associated golf club shaft into the hollow plastic tube. The flexible restricted throat element subsequently returns to its initial shape for close fitting circumferential support of the golf club hosel at an upper end of the golf club shaft adjacent the golf club head. The predetermined diameter of the hollow plastic tube is dimensioned to at least substantially peripherally engage the golf club handle at least adjacent the open lower end. The entire length of the golf club shaft is thus protected by the hollow plastic tube through the circumferential support of the golf club hosel by the flexible restricted throat opening at the open upper end and by the at least substantial peripheral engagement of the golf club handle at the lower open end so as to provide suspended non-engagement of the golf club shaft therebetween.\nA wrap for protecting a portion of a golf club shaft from abrasion within a golf bag. The inventive device includes a pad strip for circumferentially extending about a shaft of a golf club. A securing strip extends from the pad strip for securing the pad strip in an annular configuration about the shaft to protect the shaft from abrasion against an upper interior surface of a receiving tube of a golf bag.\nA golf club shaft protector is disclosed as including an elongated hollow plastic tube having a polygonal cross-sectional configuration with an unequal number of sides interconnected to each other by corner sections. Each of the unequal number of sides have the same predetermined length and each of the corner sections have the same predetermined angular shape. One of each of the corner sections faces one of each of the sides. Preferably, the unequal number of sides is at least seven to provide the largest possible opening with the greatest amount of rigidity for receiving the golf club shaft, including a golf club handle. The golf club shaft protector also includes a tubular element mounted adjacent the open upper end of the hollow plastic tube and includes an outer wall section, an inner wall section and flexible lip sections. The outer wall section surrounds an outer wall area of the tube adjacent the open upper end, the inner wall section surrounds an inner wall area of the tube adjacent the open upper end and the flexible lip sections extend over the open upper end of the tube for resiliently supporting a golf club shaft in centered position relative to the hollow plastic tube. The aforementioned elements facilitate the entry and removal of a golf club shaft including a golf club handle or grip, while protecting the golf club shaft against marring, scratching or other damage while retained within the hollow plastic tube.\nWhile these golf club shaft protection devices may be suitable for the purposes for which they were designed, they would not be as suitable for the purposes of the present invention, as hereinafter described.\nThe present invention discloses a device for protecting golf clubs from being scratched due to the golf clubs being inserted into a golf club protection tube which are commercially available. The present invention discloses a flanged tubular insert for placement into the top neck of the golf club protection tube. The present invention is padded with soft pile material extending circumferentially about the insert having a hole therein for inserting of the golf club shaft. Another embodiment is provided which has a Y-shaped upstanding member mounted on the top edge of the present invention within which is placed the golf club head.\nA primary object of the present invention is to provide a flanged insert which can be used in conjunction with existing golf club tubes for the protection of golf clubs.\nAnother object of the present invention is to provide a flanged insert which can be used in conjunction with existing golf club tubes to prevent the golf club shafts from contact with the lip of existing golf club tubes.\nYet another object of the present invention is to provide a golf club tube flanged insert having a lower portion of smaller diameter and an upper portion of greater diameter.\nStill yet another object of the present invention is to provide a golf club tube flanged insert having a soft pile material circumferentially affixed to the greater upper interior wall of said flanged insert.\nYet another object of the present invention is to provide an alternate embodiment of a golf club tube flanged insert having a soft pile material circumferentially affixed to the greater upper interior wall of said flanged insert and further having a club head bearing member comprised of a Y-shaped member fixedly attached to the exterior surface of the flanged insert therein providing means for supporting the club head within the V-portion and preventing any lateral movement of the club head which normally occurs during the course of play.\nAdditional objects of the present invention will appear as the description proceeds.\nThe present invention overcomes the shortcomings of the prior art by providing a golf club protection device comprising a flanged insert which seats by compression fitting or adhesively fixed into the top opening of existing golf club protection tubes and said flanged insert having an acrylic soft pile material circumferentially affixed to the interior wall surface of said flanged insert.\nAnother embodiment is provided having all of the properties of the preferred embodiment in addition to a club head bearing member comprised of a Y-shaped member fixedly attached to the exterior surface of the flanged insert therein providing means for supporting the club head within the V-portion and preventing any lateral movement of the club head which normally occurs during the course of play.\nThe foregoing and other objects and advantages will appear from the description to follow. In the description reference is made to the accompanying drawings, which form a part hereof, and in which is shown by way of illustration specific embodiments in which the invention may be practiced. These embodiments will be described in sufficient detail to enable those skilled in the art to practice the invention, and it is to be understood that other embodiments may be utilized and that structural changes may be made without departing from the scope of the invention. In the accompanying drawings, like reference characters designate the same or similar parts throughout the several views.\nThe following detailed description is, therefore, not to be taken in a limiting sense, and the scope of the present invention is best defined by the appended claims."}
-{"text": "As digital apparatuses provide diverse functions nowadays, digital apparatuses having a built-in camera are developed and used. Recently, there emerges a digital apparatus that automatically moves a lens of a camera built in the digital apparatus using a predetermined lens driving motor, thereby allowing a user of the digital apparatus to use the camera more conveniently.\nThe lens driving motor should be provided in a small size with consideration of miniaturization of the digital apparatus. Also, an inner portion of the lens driving motor, particularly, an electrical line through which a current flows should be stably connected, and reliability in accuracy for a movement amount of a lens should be excellent considering that the lens driving motor is an optical apparatus. Japanese Patent Publication No. 2004-280031 proposes a related art lens driving motor capable of the above-described purposed to some extent.\nThe Japanese Patent discloses a lens driving motor including a coil and a carrier coupled together, a lens coupled to the carrier, and a magnet installed on an outer side of the coil.\nThis related art lens driving motor operates in such a way that the carrier and the lens move in an upward direction using electromagnetic force generated between the coil and the magnet when a current flows through the coil. Also, one side of a spring is fixed to the carrier, and the other side of the spring is clamped by a yoke and caps of upper and lower portions so that a stop position of the carrier is indicated.\nMeanwhile, according to the Japanese Patent, lead lines of a coil through which a current is applied to the coil are soldered at an upper spring and a lower spring, respectively, to receive external power via the upper spring and the lower spring.\nHowever, the lead lines of the coil should be soldered at the upper and lower springs, respectively, during a process of assembling respective parts because of the structure of the related art lens driving motor. Since a lens driving apparatus has a very small volume of about 1.2-1.5 cm3, a soldering operation should be performed in an inside of a narrow space, which reduces workability.\nAlso, due to a yoke shape where cross-sections are bent, the lead line of the coil should be inserted into a gap between the yoke and the carrier, extracted to an upper side, and then soldered at the upper spring when the lead line is soldered at the upper spring. These processes are more difficult to perform.\nAlso, since one end of the lead line of the coil is connected to the coil, of course, and move together with the coil, but the other end of the coil is fixed to a portion (e.g., the yoke) that moves independently of the coil, friction caused by contact and an external load are applied to the lead line. These problems generate deterioration of the lead line and reduce reliability of a product."}
-{"text": "Currently available and well known plastic lined piping products comprise a family of pipes, fittings, and valves especially designed for handling corrosive or high purity liquids. Such products generally comprise steel lined with a polymeric material, e.g., polyvinylidene chloride, polypropylene, polyvinylidene fluoride, and polytetrafluoroethylene. Such products enjoy both the structural integrity of steel and the high chemical resistance characteristic of the selected polymeric liner. The external surface of such products often bears a thin coating of paint to protect the steel from corrosive substances found in the operating environment.\nEach pipe, fitting, and valve will contain a flange at or substantially near each end thereof. Adjacent pipes, fittings, and valves, within a given pipeline, may be joined one to the other by the fastening together of such flanges, e.g., by bolting. Such bolts may be coated with a polymer, e.g., polytetrafluoroethylene, to inhibit the corrosion thereof.\nTo secure the polymeric liner within a given pipe, fitting, or valve, such liner is provided with a length greater than that of the bore through which it extends. Upon the application of heat and pressure, the liner is flared over the pipe end, which in some cases will be the anterior flange face. In such cases, the flaring will preferably not cause an extension of the liner to the circumferential edge of the flange. Thus, when adjacent flanges are secured together, e.g., by bolting, a gap remains between adjacent opposing anterior flange faces equal to at least the thicknesses of the flared polymeric liners. The presence of this gap precludes the conductivity of electricity between adjacent components when non-conductive bolts, e.g. polytetrafluoroethylene coated bolts, are used.\nMethods of establishing electrical conductivity between adjacent components are known in the art. For instance, external lock tooth washers may be compressed between the bolt head and posterior flange face of a first component, and between the nut and flange of the adjacent component. Such washers bear teeth which serve to remove the paint from the outer surface of the posterior flange faces and the non-conductive coating from the bolts, thereby providing points of conductivity. While providing improved conductivity along a given pipeline, external lock tooth washers present undesirable complications. First, the installation of such washers or the replacement of a corroded washer after the assembly of the pipeline is impossible without dismantling at least the effected portion of the pipeline. This often involves clearing the pipeline of chemical substances contained therein, removing the connecting bolts, installing the washer(s), and retorquing the bolts. Second, the presence of external lock tooth washers causes torquing problems. Such washers cause an increase in friction between the nut and the posterior flange face. This increased friction is added torque, causing an inaccurate torque reading by the torque wrench. Thus, such washers may lead to insufficiently tight bolts, which may in turn lead to leaks at the pipe joint.\nIt is further known to provide conductivity clamps on adjacent components which bite into the paint and provide points of conductivity. A conductive wire connects the clamps located on adjacent pipes, providing a conductive path therebetween. Such clamps tend to be relatively expensive.\nThose in the industry would find great advantage in a means for providing electrical contact between the opposing anterior faces of adjacent metal pipe flanges which is both cost effective and practical. Such means should be easy replaceable and should not result in torquing problems. The loss or corrosion of such means should likewise be easily detectable. Such means should preferably permit the establishment of electrical conductivity between relatively widely spaced anterior flange faces, e.g., anterior flange faces spaced approximately one inch apart."}
-{"text": "    Patent document 1: JP-2003-198713 A (corresponding to US-2003/0114202)    Patent document 2: JP-2003-218996 A\nA cellular phone compliant to a Bluetooth (registered trademark) communication function becomes widely used. In connection with it, an in-vehicle handsfree apparatus for establishing a Bluetooth communication link with the above cellular phone has been offered (for example, Patent documents 1, 2).\nIn such an in-vehicle handsfree apparatus, when the cellular phone receives an incoming signal via a communication network while establishing the Bluetooth communication link between the in-vehicle handsfree apparatus and the cellular phone, the cellular phone sends an incoming notice signal or call to the handsfree apparatus. Upon receiving it, the handsfree apparatus executes the handsfree incoming control to give a notice for indicating that the cellular phone has received the incoming call while prohibiting the cellular phone from outputting an incoming ringer tone.\nIn contrast, an in-vehicle navigation apparatus having a handsfree function connectable with multiple communication terminals is offered. A cellular phone may receive an incoming signal from a communication network in a state where the in-vehicle navigation apparatus cannot execute a handsfree incoming control relative to the incoming signal. In such a case, the navigation apparatus has a difficulty in executing a handsfree incoming control. The user needs to operate the cellular phone, instead of operating the navigation apparatus, to answer the incoming signal. Herein, it is not a problem if the navigation apparatus is operated while the vehicle is parked or stopped; however, if the vehicle is running, it becomes a cause to interfere the safety of driving the vehicle. This is not desirable."}
-{"text": "A common device in the household of today is a shelf unit for holding various items for use in the bathroom, especially in a bathtub or shower area. These units have been manufactured to hold shampoos, various cream and lotion products, wash cloths, shaving products, conditioners, etc.\nOf course, the units must be readily accessible when these items are needed during the bathing process. Consequently, many units have been designed to fit over or attach to the extension arm between the wall and a shower head. A unit of this type is disclosed in U.S. Pat. No. 4,387,811 (1983).\nOther designs have been developed for these shelf units. For example, U.S. Pat. No. 5,921,410 (1999) discloses a collapsible shelf unit which can be attached to a conventional towel rack, hand rail or the like.\nHowever, the prior art units have one or more disadvantages, such as inadequate storage for certain items or lack of structural integrity. In addition, some of the \u201cshower head\u201d units require the user to reach through the shower spray to obtain the desired item.\nThus, there is a need in the industry for an effective and reliable accessory shelf unit for holding various items, especially in a bathtub or shower area."}
-{"text": "This invention relates generally to flip-down type sunglasses which are generally worn in athletic competition. The sunglasses may be moved to a horizontal positon so that they do not obscure the wearer's vision. The glasses may be quickly flipped down into a vertical position to shield the eyes from the sun. While a variety of flip-down constructions have been known in the art, such constructions have been less than completely satisfactory in that the structures required for the mounting and positioning the moveable visor have been difficult to inexpensively mass produce.\nGenerally speaking, in accordance with the invention, an improved construction for flip-down sunglasses is provided. The sunglasses have a fixed part which extends across the wearer's forehead and is attached by the use of earpieces and an elastic strap. Pivotedly mounted to the fixed part is a moveable part carrying a sun visor which may be flipped down by the wearer as desired. The moveable part has a pair of ribs which engage a groove carried on a flexible arm integral with the fixed part. The grooves engage the ribs to position the glasses either horizontally or vertically with respect to the wearer's face.\nAccordingly, it is an object of this invention to provide an improved flip-down sunglass construction that may be inexpensively manufactured.\nIt is another object of this invention to provide an improved flip-down sunglass construction that may be manufactured using molding processes.\nIt is a further object of this invention to provide an improved flip-down sunglass construction that is relatively lightweight.\nStill other objects of this invention will become apparent upon the reading of the detailed specification to follow."}
-{"text": "The present invention relates to a card connector, and more specifically, to a card connector with an ejection mechanism of a push-push type. In the card connector of the invention, it is easy to insert a card.\nPatent Reference has disclosed a conventional card connector with an ejection mechanism of a push-push type. In such a card connector, a tray is disposed to eject outside the card connector for covering a bottom surface and side surfaces of a card, so that the card is inserted into the card connector. When the card is removed from the card connector, the tray is ejected again. Accordingly, a mechanism for ejecting the tray determines how easy a card can be inserted into the card connector.\nPatent Reference: Japanese Patent Publication No. 2000-148927\nIn another type of conventional connector, instead of the tray, an ejector having a plate portion is used. In order to reduce a size of a card connector, the ejector is arranged to slide inside the card connector with a card mounted thereon, thereby not being ejected outside the card connector. Accordingly, in order to easily insert the card into the card connector, it is necessary to provide a large insertion opening for inserting the card.\nIn view of the problems described above, an object of the present invention is to provide a card connector with an ejection mechanism capable of solving the problems in the conventional card connectors.\nFurther objects will be apparent from the following description of the invention."}
-{"text": "In order to effectively utilize radio wave resources in a mobile communication system, it is required to compress speech signals at a low bit rate. On the other hand, it is expected from the user to improve quality of communication speech and implement communication services with high presence. In order to implement this, it is preferable not only to improve quality of speech signals, but also to be capable of encoding signals other than speech, such as audio signals having a wider band with high quality.\nFurther, in an environment where various types of networks are present, a speech coding scheme is required that can flexibly support communication between different networks, communication between terminals utilizing different services, communication between terminals having different processing performance, and conversational communication at multipoints as well as communication between two parties.\nMoreover, a speech coding scheme is required to be robust against transmission path errors (in particular, packet loss in packet switching networks typified by IP networks).\nOne speech coding scheme satisfying such requirements is the bandwidth scalable speech coding scheme. The bandwidth scalable coding scheme is a coding scheme that encodes speech signals in a layered way, and a coding scheme where coding quality increases in accordance with an increase in the number of coding layers. The bit rate can be set variable by increasing or decreasing the number of coding layers, so that it is possible to effectively use transmission path capacity.\nFurther, with the bandwidth scalable speech coding scheme, it is only necessary to receive at least the data coded by a base layer at a decoder side, and it is possible to allow to some extent information coded by additional layers being lost on the transmission path, and therefore the bandwidth scalable speech coding scheme provides robustness against transmission path errors. Further, the frequency bandwidth of speech signals to be encoded also becomes wider in accordance with an increase in the number of coding layers. For example, for a base layer (i.e. core layer), a coding scheme for telephone band speech of the related art is used. Further, in additional layers (i.e. enhancement layers), layers are configured so that wideband speech which has a bandwidth such as 7 kHz can be encoded.\nIn this way, with the band scalable speech coding scheme, telephone band speech signals are encoded in the core layer, and high-quality wideband signals are encoded in the enhancement layers, so that it is possible to utilize the bandwidth scalable speech coding scheme for both telephone band speech service terminals and high-quality wideband speech service terminals and support multipoint communication including the two kinds of terminals. Further, the coded information is layered, so that it is possible to increase error robustness by devising a transmission method, and readily control the bit rate on the encoding side or on the transmission path. Therefore, the bandwidth scalable speech coding scheme draws attention as a speech coding scheme for future communication.\nThe method disclosed in non-patent document 1 is given as an example of the bandwidth scalable speech coding scheme described above.\nIn the bandwidth scalable speech coding scheme disclosed in non-patent document 1, MDCT coefficients are encoded using a scale factor and fine structure information for each band. The scale factor is Huffman encoded, and the fine structure is subjected to vector quantization. Auditory weighting of each band is calculated using a scale factor decoding result, and the bit allocation to each band is decided. The bandwidth of each band is non-uniform and set in advance so as to become wider for a higher band.\nFurther, transmission information is classified into four groups as described below.\nA: Core codec coding information\nB: High-band scale factor coding information\nC: Low-band scale factor coding information\nD: Spectrum fine structure coding information\nFurther, the following processing is carried out on the decoding side.\n<Case 1> When information for A cannot be received completely, decoded speech is generated by carrying out frame erasure concealment processing.\n<Case 2> When only information for A is received, a decoded signal for the core codec is outputted.\n<Case 3> When information for B is received in addition to the information for A, a high band is generated by mirroring the decoded signal for the core codec and a decoded signal having a wider bandwidth than the decoded signal of the core codec is generated. Decoded information for B is used in generation of high band spectrum shapes. Mirroring is carried out at a voiced frame, and is carried out so that the harmonic structure does not collapse. The high band is generated at an unvoiced frame using random noise.<Case 4> When information for C is received in addition to information for A and B, the same decoding processing as in case 3 is carried out using only information for A and B.<Case 5> When information for D is received in addition to the information for A, B and C, complete decoding processing is carried out at bands where all information for A to D is received, and a fine spectrum is decoded by mirroring a decoded signal spectrum on the low band side at bands where information for D is not received. Even if the information for D is not received, it is possible to receive the information for B and C, and this information for B and C is utilized in decoding of spectrum envelope information. Mirroring is carried out at a voiced frame, and is carried out so that the harmonic structure does not collapse. The high band is generated at an unvoiced frame using random noise.Non-patent document 1: B. Kovesi et al, \u201cA scalable speech and audio coding scheme with continuous bit rate flexibility,\u201d in proc. IEEE ICASSP2004, pp. I-273--I-276."}
-{"text": "Automated keyword generation provides valuable advantages when producing searchable content, such as advertising and media content. When a user is browsing on the internet, the usual method by which a user reaches a website is through a search engine, and one factor in how prominently the search engine displays the website is the quality of keywords associated with the website. For search engine marketing (SEM) campaigns, bidding on the right keywords is extremely important. While keyword suggestion tools currently exist, these existing tools often produce overly broad or inapplicable keyword results. Further, in the case that a content provider is generating large amounts of searchable content, any amount of time refining keyword results from existing keyword suggestion tools quickly becomes a major hurdle. Specifically, within the context of generating automated advertising campaigns for quickly changing sources of content, it becomes extremely advantageous for a user to minimize time required to focus and refine keyword generation results to any content produced or content of interest."}
-{"text": "Computer systems generally include one or more processors interfaced to a temporary data-storage device such as a memory device and one or more persistent data-storage devices such as disk drives. Each disk drive generally has an associated disk controller. Data is transferred between the disk drives and the disk controllers. Data is also transferred between the disk controller(s) and the memory device over a communications bus or similar.\nData organization in a computer system such as that above is important in relational database systems that deal with complex queries against large volumes of data. Relational database systems allow data to be stored in tables that are organized as both a set of columns and a set of rows. Standard commands are used to define the columns and rows of tables and data is subsequently entered in accordance with the defined structure.\nThe defined table structure is locally maintained but may not correspond to the physical organization of the data. In a parallel shared nothing relational database data can be stored across multiple data-storage facilities, each data-storage facility in turn including one or more disk drives. Data partitioning can be performed in order to enhance parallel processing across multiple data-storage facilities. The intent behind data partitioning is to evenly distribute the data among all computational elements such that performance scales linearly as more computational elements are added. The database could be hash partitioned, range partitioned, round-robin partitioned or not partitioned at all.\nHash partitioning is a partitioning scheme in which a predefined hash function and map is used to assign rows in a table to respective processing modules and data-storage facilities. The hashing function generates a hash or partition bucket number and the partition numbers are mapped to data-storage facilities. Range partitioning is a partitioning scheme in which each data-storage facility manages the records falling within a range of values. Round Robin partitioned is a partitioning scheme where the data-storage facility is picked in a round robin fashion. No partitioning means that a single data-storage facility manages all of the rows.\nOne drawback of current systems is that the mapping of partitions to the data-storage facilities on which the rows are stored is often required to be physically defined in advance especially when the mapping is not to all the data-storage facilities but to a subset of them. The mapping of partitions to data-storage facilities is therefore static. A static mapping to specific user defined data-storage facilities often leads to uneven distribution of rows over the data-storage facilities. This in turn has the potential to increase the execution time for complex queries.\nUsers of relational database systems require the minimum time possible for execution of complex queries against large amounts of data. In a parallel shared nothing relational database system it is often important to evenly allocate both table rows and free space across multiple data-storage facilities."}
-{"text": "Medical instruments often need to be temporarily connected to peripheral devices and components in the course of operation. An example may be a sampling tube connected to an analyzing instrument such as a capnograph. Another example may an ultrasound probe connected to a sonographic imaging instrument. Such peripheral devices may need to be replaced frequently due to one or more reasons. For example, a disposable probe may be used for each treated patient, and should be replaced after use by a new probe for a next patient. Another reason for frequently connecting and disconnecting probes from an instrument may be related to multi-purpose instruments. Such instruments are configured to carry our one of several routines, for obtaining one of several optional purposes. Generally, a particular routine and purpose may be associated with a specific peripheral device that needs to be connected to the instrument for carrying out the routine. For example a user of a sonographic imaging system may wish to carry out one type of examination using one type of ultrasound probe, and then to carry out a second type of examination using a second type of ultrasound probe. Thus, frequent replacement of probes is required, typically being carried out by disconnecting a previously used probe and connecting a new probe to the instrument, instead."}
-{"text": "Due to the ever increasing demand for data storage in today's digital world, distributed storage systems have received a lot of attention. In distributed storage systems, data as a whole or different portions of data are stored over several physically independent computer storage servers (also known as storage nodes). The total storage capacity of the distributed storage system is designed to be much greater than that of a single storage node.\nThere is an exponentially increasing amount of data produced and stored worldwide in distributed storage systems. A large fraction of that data is duplicated, for instance, as a result of unintentional copies of all or part of the data being stored multiple times, therefore leading to wasted storage resources. In particular, company/enterprise data, consisting mostly of structured formats such as documents and spreadsheets, is often replete with unnecessary duplicate data sometimes estimated over \u02dc70%. Therefore, an enterprise-wide storage system can make significant space savings if the data can be de-duplicated and stored only once. However, in many distributed database systems, deduplication is difficult to implement as the data models supported by, for example, relational database management systems (RDBMS), are not designed with deduplication in mind.\nFile-system level deduplication products are alternatives that aim to work with raw data on disks underneath any higher level storage abstractions such as databases, thereby leaving the upper or higher layers to operate as normal (e.g., unaware of the deduplication process underneath it). However, it has been shown that file-system level deduplication does not lead to optimal storage resource savings (e.g., for databases, especially) as there are too many small changes (e.g., file headers, timestamps, etc.) that make blocks of data subtly different and thereby defeat deduplication efforts. File-system level deduplication can also introduce additional data loss risks, as the database on top can make incorrect assumptions about which critical data could be lost due to a physical disk error.\nIn view of the above, there is a need for efficient systems and methods for data storage that can overcome these challenges and store the data in a distributed storage system securely, reliably and with minimal duplication."}
-{"text": "The Internet Engineering Task Force (IETF) has defined protocols in the area of Layer Two Tunneling Protocol\u2014Version 3 (L2TPv3) and Pseudo Wire Emulation Edge-to-Edge (PWE3) Architecture, as defined in RFC 3193, RFC 3931, RFC 3985, RFC 4591, and RFC 4719. The L2TPv3 pseudo wires may support multiple types of layer two encapsulations or \u201cPseudowire Types\u201d, including protocols similar in packet format to IEEE 802.3, virtual LAN tagging protocols such as IEEE 802.1Q, IEEE 802.1ad, and IEEE 802.1ah, as well as Point-to-Point Protocol (PPP), Frame Relay, and other layer 2 technologies.\nAs defined in the standard track protocols above, the Pseudowire Types including Ethernet, VLAN, PPP, or Frame Relay, etc are encapsulated into \u201cL2TPv3\u201d to emulate a single layer 2 tunnel link between exactly two endpoints.\nEnabling L2TPv3 on cable modems (CM), embedded multimedia terminal adapters (E-MTA), cable modem routers (CM-R), cable modem with Circuit Emulation Service (CES), cable modem with IPSec or Firewall and/or similar DOCSIS Customer Premise Equipment (D-CPE) will allow L2TPv3 IP packets to be encapsulated into Data Over Cable Service Interface Specification (DOCSIS) data frames for two-way data transmission in a cable system. Devices with this capability are referred to herein as \u201cL2TPv3 enabled DOCSIS CPE\u201d, or simply \u201cEnhanced CPE\u201d."}
-{"text": "The transport of truck trailers by train has traditionally been limited to long haul transport. Trailer design, railroad car design and scheduling problems have made it difficult for trains to compete in short and medium haul corridors. To date, trailers have had to be reinforced to handle the stress of rail transport. In addition, trucking companies have been forced to accept fairly wide windows of time for delivery.\nCurrently, trucking customers arrive at the train terminal and wait in line for a clerk to manually input the details of the shipment at the gate, to include billing information, contents of the truck, owner, driver, destination, weight, commodity code, commodity description, package type, broker and other pertinent information needed to transport the trailer and provide billing information. This process is time consuming and tedious; tractor-trailer drivers may wait in line at the gate for over an hour as each truck ahead of his goes through the check-in process. In addition, the time required to check-in a trailer is lengthy.\nFurthermore, transport of trailers via rail has been unpredictable. When the trucking customer arrives with a trailer to transport, they do not know if their trailer will be transported on the next train. The customer's container may be left behind for subsequent trains. In addition, it can be difficult to determine the train the trailer will be on and when that train will arrive at its destination. As a result trailer transportation via rail is an unreliable, frustrating and complex system.\nTrucks are a very efficient transportation medium, especially in short and medium haul corridors. Over the last 10 to 15 years, there have been a variety of attempts by the railway industry to move truck trailers onto rails. As noted above, most of these attempts to create intermodal transportation systems, however, require expensive modifications to the trailer to make it sturdy enough to withstand forces applied by cranes and by transport on the rails. Also the cost of building and outfitting railway terminals for truck transport is very expensive. A complex railway system for loading and loading trains is needed. In addition, trucking companies chafe at the long delays and complex processes of the railway companies.\nWhat is needed is a simple and reliable transport system which ensures that a train leaves on schedule and arrives at its destination within a predetermined, narrow window of time.\nIn addition, what is needed is a low cost high quality profitable intermodal system that reduces current bottlenecks and provides an efficient means of transporting trailers via rail.\nFinally, what is needed is a system and method for efficiently and economically transporting trailers via train across high density corridors."}
-{"text": "1. Technical Field\nThe present invention is related to an apparatus and method for moving a first object relative to a second object and holding them in a defined position. In particular, the present invention relates to a scanning system for micromechanical devices.\n2. Description of the Prior Art\nWhen an object is to be moved or driven with a defined precision, then a device that controls the object's position and movement is employed. Such a device is a scanner that includes an actuator, whereby various kinds of driving principles are known for those actuators. Well known are electromagnetic, electrostatic, superconducting, piezoelectric, ultrasonic, pneumatic, air-based, thermal expansion, bimetal, and shape-memory alloy actuators. Micro-actuators and the micro machines driven by them are expected to have wide varieties of applications in information technologies, in medical, bioengineering, optics, and robotics fields. For example, electromagnetic actuators are ubiquitous. They can be found in everything from a large train to the smallest time piece. However, even the smallest magnetic actuators are usually made with wound coils and magnetic materials.\nIn the European patent application publication no. EP 0 998 019 A1, currently owned by the present applicant, a magnetic scanning or positioning system with at least two degrees of freedom is described. The magnetic scanning or positioning system comprises a supporting base equipped with magnets, a movable platform equipped with at least two electrical coils, and suspension elements providing an elastic connection between the movable platform and the supporting base. The electrical coils are positioned flat on the movable platform, thereby forming an essentially flat arrangement with the movable platform. The scanning or positioning with at least two degrees of freedom can be used in the field of scanning probe microscopy or in the field of data storage or imaging.\nAnother example is known as an electrostatic actuator. European patent application publication no. EP 0 865 151 A2 describes such an electrostatic actuator. The electrostatic actuator comprises a first member and a second member. The first member has a first opposed surface that includes an array of driven electrodes with pitch and the second member has a second opposed surface including an array of drive electrodes. A support positions the first member adjacent the second member with the first opposed surface spaced apart from the second opposed surface by a spacing. The ratio of the spacing and the pitch should be less than eight, and is preferably less than 2.25. The support permits the first member to move relative to the second member, or vice versa. A voltage source establishes a spatially substantially alternating voltage pattern on the array of driven electrodes. An electrode control establishes a substantially alternating voltage pattern on the array of drive electrode, and selectively imposes a local disruption on the substantially alternating voltage pattern on the array of drive electrodes to move the movable one of the first member and the second member relative to the other.\nEuropean patent EP 0 071 066 B1, granted to the present applicant, discloses an electric traveling support. The traveling support comprises a piezoelectric plate resting on three legs whose bottom surface is insulated from the bench on which the support is to travel, by a dielectric. The piezoelectric plate can be caused to contract by means of an actuating voltage applied via sliders to top and bottom electrodes on the plate. By applying a voltage to the legs, these may be clamped selectively by electrostatic forces effective across the dielectric. Appropriate control of the actuating and clamping voltages causes the support to either move in a linear or rotary fashion.\nAll the aforementioned actuators and scanners need energy not only for movement but also for holding and positioning. Low power precise positioning of parts and devices become more and more important as (nano)micro-mechanics gains in importance. In particular in nano-mechanics on the bases of local probes where tips have to be positioned with nano-meter precision reliable scanners are missing. For piezo-scanners domain wall creep and for electromagnetic scanners thermal creep are unsolved problems. Additionally, present scanners consume energy for holding and restoring the last position. Furthermore for conventional scanners the energy needed for moving from position to position increases with increasing spring deformation.\nIt is therefore an object of the present invention to overcome the disadvantages of the prior art.\nIt is another object of the present invention to provide a scanner with integrated tracking and nearly zero power consumption for holding a position with high precision."}
-{"text": "Crawling, flying, and pestilent insects are troublesome to catch and dispose of, particularly for individuals who are afraid of, or repulsed by, such creatures. Existing devices such as fly swatters are currently used to crush such insects. However, the messy carcass of the dead insect is left remaining on the walls and/or on the fly swatter. Prior art fly swatters are known to have used frontal adhesives, but the frontal planes of the fly swatter do not provide a readily available means of rapidly, effectively, and hygienically entombing and completely sealing the flying and crawling insects. Moreover, such devices also fail to adequately seal the exposed adhesive surfaces on the device after insect contact, adhesion, and sealing. This is necessary to prevent future coincidental adhesive contact with surrounding surfaces or objects, including the user of the device.\nAnother disadvantage of prior art devices is that insect contact is established on an essentially random basis, versus the will of an operator on an immediate need basis. Additionally, use of chemical insecticides in human habitats may be undesirable and may be harmful, either short or long term. This is particularly true in nurseries, and/or in the presence of the very ill or elderly. The present invention eliminates the need for chemical insecticides in human habitats where their use could be potentially harmful.\nAdhesive technology exists today wherein the adhesive qualities can be modified to specifically address the user's requirements or desires. Additionally, pressure-sensitive adhesive papers or other materials are manufactured today at low unit-costs. Such papers might include the incorporation of complex or simple die-cuts and folds. Hence, what is needed is a device which utilizes such adhesive technology providing an operator means for selectively capturing and quickly disposing of an insect.\nStill another disadvantage of the prior art which uses adhesive is that stationary insect traps such as fly paper can accumulate large numbers of insects. Such accumulated pests might carry diseases or pathogens which could present a health-hazard potential. These accumulated insects, either dead or alive, are exposed to human surroundings for possible air borne or physical contact by humans and pets.\nA solution is needed which affords an operator a willful, portable means for quickly capturing and disposing of an insect in a hygienic and clean manner, even when an insect is at arms reach. This device can thereby incorporate an extension handle attached to a holder for an adhesive paper mounting device. Alternatively, a fixed location unit could also be provided which uses the same hygienic capturing and disposal system. The resulting portable or fixed device should allow for totally containing, sealing, and disposing of an insect in a fast, convenient and hygienic manner."}
-{"text": "1. Technical Field\nThis invention generally relates to the field of projectiles, and more specifically relates to projectiles with deployable members.\n2. Background Art\nMany advances have been made in the art of projectiles, such as bullets fired from guns. Several known bullets are made of lead or other soft material that expands (known as \u201cmushrooming\u201d) when the bullet hits. The expansion of a lead bullet inside a target causes a greater knock-down effect, and increases the damage done to bones and internal organs, but typically slows the bullet to the point that it does not exit the target. In hunting applications, it is desirable for the bullet to exit the animal so the animal bleeds from the exit would, allowing the hunter to track the animal from the trail of blood. One way to assure the bullet exits the animal is to use a harder material that does not expand upon impact. The drawback of this approach is the damage done to the animal is not as great as for a softer, expanding bullet, increasing the likelihood of survival for an animal shot with a hard bullet. In addition, because a hard bullet does not expand, the animal will not likely bleed a great deal because the exit would is small, the same diameter of the bullet.\nSome projectiles have been developed with members that deploy to increase the damage when the projectile hits its target. For example, U.S. Pat. No. 6,240,849 to Holler and U.S. Pat. No. 1,464,032 to Daynix disclose projectiles that have members that deploy in-flight. These members increase the damage to the target upon impact. U.S. Pat. No. 1,318,858 to Frick discloses a projectile that may expand in-flight, or that may expand upon impact with a target. The Frick projectile includes pivoting knife arms that extend to create more damage to the target. The configuration of the Frick projectile is quite complex, and would be very difficult to manufacture in a cost-effective manner. What is needed is a projectile that provides members that deploy upon impact with a target that may be manufactured and assembled in a cost-effective manner."}
-{"text": "In the standardization of 3GPP (Third Generation Partnership Project) there is on going work to standardize user authentication routines, especially for the so called Generic Authentication Architecture (GAA) involving a mutual authentication between a client and an application server. In the communication network several different applications will be available for the client and these applications will be supplied by third party suppliers, i.e. different from the Service Provider (SP). However, a client accessing several different applications would like to make only one single authentication, a so called Single Sign On (SSO) service, providing the possibility for the user to only authenticate once during a session, rather than to authenticate to each new application server it wants to use. This will make the authentication procedure much easier for the client. Also this kind of authentication service can be offered to third party application service providers as a service from the Service Provider handling the network. The GAA is aiming for solving this problem and making such a service available in the 3GPP network. The generic authentication provides an authentication of the users on an application level based on the proven security mechanism in the Public Land Mobile Network (PLMN).\nThe GAA is specified through the 3GPP group and drafts of the specification may be obtained through their web site, for instance the documents TS 32.220 and TR 33.919 maybe mentioned as good starting points regarding GAA. The GAA system may be explained as follows: A number of applications share a need for mutual authentication between a client/user (called UE, User Equipment, in the standard) and an Application Server (AS) in order to allow for further communication. This is necessary when the user wants to access servers demanding authentication, e.g. content servers charging for their services, certificate demanding web sites (e.g. banks), and similar application servers. Since many applications share the same need, it has been considered to specify a Generic Authentication Architecture (GAA), providing the architecture for allowing application servers access to the infrastructure authentication systems. Thus, if the application server trusts the service provider, this architecture may simplify the authentication schemes both for the user and for the application servers. The user needs only to authenticate once during a session, rather than authenticating towards every application server accessed.\nIn a GAA based session the user authenticates with the network infrastructure by providing an ID to a Bootstrap Function (BSF), this ID may for instance consist of the UE IMSI number (International Mobile Subscriber Identifier), which is a unique number coupled identifying a user. The IMSI is relayed to the Home Subscriber System (HSS, also called Home location register system) and the HSS provides an Authentication Vector (AV) to the BSF. The BSF authenticates the UE based on the USIM (Universal Subscriber Identity Module) and the UMTS-AKA methods, and sends a TID (Transaction Identifier) to the UE.\nAlso at the same time, work is on going to standardize so called Flow Based Charging (FBC). FBC has the aim to make it possible to charge users for service usage on a finer granularity than what is possible today. For instance it is of interest to identify the type of IP session that a user is running, the type of applications involved and so on. For instance one would like to be able to differentiate the charging costs for different types of services, e.g. streaming video may be charged more than exchanging plain text messages like simple email messages. There are many different services that may be used including both user to user and user to network services. Service data flows from these services may be identified and charged in many different ways. The FBC method is used to set up charging filters that is used by the CRF (Charging Rules Function) for different applications.\nThe filters provided in the FBC method can be quite complex and may involve source and destination address, source and destination port number, and transaction protocol, enabling a fine granularity of charging. Charging models requiring even more complex data may use special filters that look further into data packets and may be defined by the TPF (Traffic Plane Function) and invoked by the CRF.\nHowever, the above-mentioned two standardization works are not today aiming for an integration of the services they provide. This will in the future be crucial in order to be able to provide differentiated charging rules depending on the user and application connected to and at the same time ensuring the authenticity of the user towards the specific application server involved in the transaction. It should not be possible to, by only providing a fake IP number, getting access to services intended for other users or obtaining services at a wrong charging rate.\nThe work in 3GPP on GAA and FBC has been so far performed in parallel, and there is no concern taken to reuse functionality between the two functions. At some stage in standardization, interoperability between GAA and FBC needs to be built in. An integrated architecture will be necessary.\nSpecifically, the problem of supporting user specific charging for GAA authenticated users is not solved in FBC. The problem may be exemplified as follows:                1. A service provider has an application server providing a service, e.g. downloading a music file. This service is accessible for mobile users via the GPRS access network on the Gi interface. GAA ensures that the user is authenticated.        2. When the user downloads a music clip, he is charged by volume because the GGSN performs byte counting on the packets downloaded, and reports CDRs to the billing system. FBC can be used to enable that a specific rating is applied to the downloaded music file. FBC provides filtering in the GGSN that counts the specified service individually.        3. Now, the service provider might at some point provide the service for a special price for some users. In that case, it is needed to provide a specific filter for individual users. How that can be done is not specified today.        \nIn another standardization implementation, work is on going regarding policy decision functions such as the PDF (Policy Decision Function) for policy control of IP bearer resources, such as Quality of Service (QoS) for a specific user. Policy Decision Function relates a level of Quality of Service to a specific user and instant in time, enabling for instance a better Quality of Service for a certain type of application such as streaming video applications, or for a customer prepared to pay more for a higher quality of connectivity. It would be efficient to incorporate the PDF function with the GAA function in order to implement a more efficient network system."}
-{"text": "The development of a method for decomposing or removing nitrogen oxides or other harmful pollutants in the exhaust gas from automotive vehicles has been increasingly demanded due to environmental pollution. A three way catalytic converter, for removing carbon monoxide, hydrocarbons and nitrogen oxides concurrently, is a practical catalyst for purification of the exhaust gas. Cordierite coated with gamma-alumina or other heat resistant carrier supporting Pd, Pt, Rh or other noble metal is for general use.\nHowever, in a lean-burn engine and a diesel engine, which are operated under the condition of high oxygen concentration the three way catalytic converter is ineffective because the noble metals are poisoned with excess oxygen. Although zeolites have been studied as the catalyst serving under the condition of high oxygen concentration, their structure easily breaks under the hydrothermal conditions. Oxide catalysts superior in heat resistance have also been studied, however, it is insufficient in catalytic activity for their practical use.\nOn the other hand, in thermal power plants or other fixed combustion equipments, nitrogen oxides in the exhaust gas containing excess oxygen are reduced by the oxide catalyst with ammonia as a reductant. Because ammonia has toxicity, a safety problem is created when this system is utilized for automotive vehicles or other moving combustion equipments or the fixed combustion equipments installed in an urban area.\nRecently, a method of removing nitrogen oxides in the exhaust gas containing excess oxygen was developed. In the method, nitrogen oxides are absorbed by an absorbent in the lean-burn region or in the presence of excess oxygen, and nitrogen oxides are desorbed and reduced by means of the three way catalyst in the rich-burn region or in the vicinity of the stoichiometric air/fuel ratio.\nKnown absorbents of nitrogen oxides are complex oxides of the Ba--Cu--O system and the Mn--Zr--O system. However, a material having better and stable absorbability even at high temperatures is demanded."}
-{"text": "The present invention relates to a motor-generator set and more particularly to a three-phase constant frequency motor-generator having only one single rotor for both the motor and the generator, the output frequency of the generator being constant notwithstanding the fact that the rotational speed of the rotor may vary."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to the field of corn breeding. In particular, the invention relates to corn seed and plants of the variety designated CV181138, and derivatives and tissue cultures thereof.\n2. Description of Related Art\nThe goal of field crop breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits include greater yield, better stalks, better roots, resistance to insecticides, herbicides, pests, and disease, tolerance to heat and drought, reduced time to crop maturity, better agronomic quality, higher nutritional value, and uniformity in germination times, stand establishment, growth rate, maturity, and fruit size.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant. A plant cross-pollinates if pollen comes to it from a flower on a different plant.\nCorn plants (Zea mays L.) can be bred by both self-pollination and cross-pollination. Both types of pollination involve the corn plant's flowers. Corn has separate male and female flowers on the same plant, located on the tassel and the ear, respectively. Natural pollination occurs in corn when wind blows pollen from the tassels to the silks that protrude from the tops of the ear shoot.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform corn plant hybrids requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods used to develop inbred plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more inbred plants or various other broad-based sources into breeding pools from which new inbred plants are developed by selfing and selection of desired phenotypes. The new inbreds are crossed with other inbred plants and the hybrids from these crosses are evaluated to determine which of those have commercial potential.\nNorth American farmers plant tens of millions of acres of corn at the present time and there are extensive national and international commercial corn breeding programs. A continuing goal of these corn breeding programs is to develop corn hybrids that are based on stable inbred plants and have one or more desirable characteristics. To accomplish this goal, the corn breeder must select and develop superior inbred parental plants."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to a boot for downhill skiing.\nIn a generally known manner, such downhill boots have a shell base and an upper.\nThe present invention relates more particularly to a rear entry ski boot, i.e., a boot having a rear portion referred to as a rear spoiler which is journaled on the shell base around a transverse axis. The front portion of the upper can be formed out of a single piece with the shell base or be a separate piece referred to as a cuff which is journaled on the shell base around a transverse axis.\nThe cuff and the rear spoiler are, depending upon the model, generally journaled around the same transverse axis or around two different axes, while for certain boot models, the rear spoiler is journaled on a transverse axis of the cuff.\nThe present invention relates to rear entry ski boots which comprise a rear spoiler journaled on the shell base around a transverse axis, regardless of the more specific arrangements which have been just described.\nThe present invention relates in a more specific manner to those boots which comprise an adjustment apparatus activated by a control element which must be tensioned during closure of the boot before skiing.\nAn adjustment apparatus of the aforementioned type can, for example, be adapted to assure the maintenance of the foot or of the lower leg of the skier or be adapted to limit the flexion of the upper of the boot.\nIt has previously been proposed to link the control element of such an apparatus to the rear spoiler of the boot such that it be tensioned by the movement of the rear spoiler with respect to the shell base during closure of the boot. These systems have the disadvantage of a relatively complex construction and are generally raised in relief with respect to the general profile or shape of the boot, which should be avoided."}
-{"text": "My invention relates to shaving mugs. The primary objectives are to eliminate the adverse characteristics of the currently existing mugs and provide a person with a reuseable shaving mug designed specifically to hold a typical shaving soap cake in a manner that is truly convenient to the user, even in travel environments. The adverse characteristics of the currently existing shaving mugs are:\na. The soap cake easily dislodges allowing the soap cake to flop around during use or to fall into the sink during clean-up. PA1 b. A lid is not provided or is not secured, allowing other items in luggage to become soiled when traveling. PA1 c. Breaks easily if dropped during use or luggage handling. PA1 a. Prevent the soap cake from dislodging during use, . . . cleaning, . . . handling, . . . or storage. PA1 b. Provide a space to contain the lather during use. PA1 c. Permit the soap cake to dry quickly between use. PA1 d. Protect the soap cake from contamination between use, especially during travel. PA1 e. Prevent soap and water from soiling clothing and other items in luggage during travel. PA1 f. Not break when dropped from normal sink heights or during typical rough handling of luggage."}
-{"text": "The present disclosure relates to a measurement system for the measurement of dry powder based agents.\nIn order to certify a dry powder fire suppression system onboard a vehicle such as an aircraft, the agent is discharged into the protected volume and an analyzer simultaneously records the amount of fire extinguishing agent in various zones of the protected volume. The amount of agent must be above some predetermined level which has been established sufficient to extinguish all possible fires for some period of time simultaneous in all zones.\nThe analyzer must be calibrated and traceable such that analyzer output proves the dry powder fire suppression system is capable of extinguishing any fire within the protected space. No known systems are capable of both measuring aerosol cloud fire extinguishing agent concentrations and being calibrated so as to measure the agent concentration for an aircraft dry powder fire suppression system certification test."}
-{"text": "1. Field of the Invention\nThis invention relates generally to a fiber modified with carboxyl groups and, more specifically, to a cellulose or acrylonitrile fiber to which methacrylic acid or a hydroxyalkyl methacrylate is graft-copolymerized. The present invention is also directed to a process for producing such a modified fiber.\n2. The Prior Art\nOne known method of introducing carboxyl groups into a cellulose fiber includes reacting monochloroacetic acid with the cellulose fiber. Because of the introduction of carboxylic groups, the cellulose fiber thus obtained has improved hydrophilicity and is capable of absorbing basic, bad odor substances such as amines and ammonia. The known method has, however, a problem that it is difficult to homogeneously react the monochloroacetic acid with a solid, bulky mass of cellulose fibers because the reaction is exothermic.\nIt is also known to introduce ether linkages into a cellulose fiber by reaction with a vinyl compound according to the Michael reaction. With this method, however, it is not possible to effectively introduce carboxyl groups into the cellulose fiber when an unsaturated carboxylic acid such as acrylic acid or maleic acid is used as the vinyl compound. Thus, in order to introduce carboxylic groups into the cellulose fiber by utilizing the Michael reaction, it is necessary to first react the cellulose fiber with acrylonitrile and to hydrolyze the cyano groups into carboxyl groups.\nAn acrylonitrile fiber has properties similar to wool but has a defect that the moisture or damp absorbing power is poor. No effective method is however known to impart hydrophilicity to the acrylonitrile fiber.\nA method is known in which a vinyl monomer having an acidic group is graft-copolymerized to a polyester or polyamide fiber using am organic peroxide such as benzoyl peroxide."}
-{"text": "This application is an application under 35 U.S.C. Section 371 of International Application Number PCT/FR991/01691 filed on Jul. 9, 1999.\nA subject matter of the present invention is a novel terephthalic polyester composition, its process of preparation by transesterification/condensation of poly(ethylene terephthalate) and of polyethylene glycol, and its use as soil-release agent in detergent formulations, in rinsing, softening or finishing formulations, for the washing, with or without pretreatment, the rinsing, the softening or the finishing of textiles, in particular polyester-based textiles.\nThe soil-release activity of ethylene terephthalate/polyethylene oxide terephthalate copolymers in the finishing of textiles, in particular polyester-based textiles, and the use of said copolymers as soil-release and/or antiredeposition agents in detergent formulations for the washing, with or without pretreatment, of textiles, in particular polyester-based textiles, are well known (U.S. Pat. Nos. 396,252; 4,116,885; 4,785,060).\nThese copolymers can, for example, derive from the transesterification/condensation of poly(ethylene terephthalate) and of polyethylene glycol (U.S. Pat. No. 4,785,060).\nThe Applicant Company has found a novel terephthalic polyester composition which can be obtained by transesterification/condensation of a poly(ethylene terephthalate) and of a polyethylene glycol and which in particular exhibits particularly good soil-release properties.\nAccording to the invention, it is a terephthalic polyester composition (TPC) comprising, as a mixture\nan ethylene terephthalate homooligomer (PET1) essentially comprising oxyethylene terephthalate (TE) repeat units of formula (I)\nxe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94Oxe2x80x94xe2x80x83xe2x80x83(I)\nxe2x80x83where A represents the 1,4-phenylene group,\nand a block terephthalic copolymer (PET2/TE-POE) comprising\nat least one polyethylene terephthalate block (PET2) composed of oxyethylene terephthalate (TE) repeat units of formula (I)\nxe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94Oxe2x80x94xe2x80x83xe2x80x83(I)\nxe2x80x83where A represents the 1,4-phenylene group,\nand at least one polyoxyethylene terephthalate block (TE-POE) of formula\nxe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94Oxe2x80x94(CH2xe2x80x94CH2xe2x80x94O)nxe2x80x94\nxe2x80x83the value of n being such that said block exhibits a number-average molecular mass of the order of 1500 to 4000, preferably of the order of 3000 to 4000,\nsaid composition being characterized in that:\nthe amount of (TE) units of the polyethylene terephthalate (PET1) does not represent more than 10%, preferably not more than 7%, of all the (TE) units present in the terephthalic polyester composition (TPC),\nthe amount by weight of all the (TE) units present in said (TPC) composition represents at least 11%, preferably from 11.5 to 17%, of the weight of said (TPC) composition,\nthe amount by weight of mono(oxyethyleneoxy) (OEO) residues of formula\n1/2Oxe2x80x94CH2xe2x80x94CH2xe2x80x94O1/2\nxe2x80x83represents at least 1.3%, preferably from 1.3 to 2.3%, of the weight of said terephthalic polyester composition (TPC),\nsaid (OEO) residues belonging to the oxyethylene aromatic diester (OAD) groups of formula\nxe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94Oxe2x80x94C (O)xe2x80x94Axe2x80x94\nxe2x80x83present in all the (PET1) and (PET2) blocks,\nand in that the weight-average molar mass of said block terephthalic copolymer (PET2/TE-POE) is at least 30000, preferably at least 35000, very particularly at least 40000.\nThe ends of the chains of (PET1) homooligomers are generally composed of xe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94OH units.\nThe ends of the chains of block terephthalic copolymer (PET2/TE-POE) are generally composed of\nxe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94OH\nunits and/or\nxe2x80x94(CH2xe2x80x94CH2xe2x80x94O)nxe2x88x921xe2x80x94CH2xe2x80x94CH2xe2x80x94OH\nunits.\nThe characteristics of the (TPC) composition of the invention can be determined by subjecting said composition to the following analyses:\nsize exclusion chromatography (SEC) using a size exclusion chromatography device with simultaneous detection by refractometry and ultraviolet, using N,N-dimethylacetamide (DMAc) comprising 10xe2x88x922 mol/l of LiBr as eluent at 100xc2x0 C.\nThis measurement (UV chromatography) makes it possible:\nto detect the presence of the A (terephthalic) aromatic nuclei,\nto determine the percentage of ethylene terephthalate homooligomer (PET1) in the terephthalic polyester composition (TPC),\nto determine the weight-average molar mass of the block terephthalic copolymer (PET2/TE-POE), expressed in polystyrene equivalents.\nproton nuclear magnetic resonance (1H NMR) using an NMR spectrometer in a CDCl3/deuterated hexafluoroisopropanol/deuterated trifluoroacetic anhydride mixture.\nThis measurement makes it possible to detect\nthe aromatic (terephthalic T) nuclei,\nthe oxyethylene units,\nthe mono(oxyethyleneoxy) (OEO) residues present in the oxyethylene aromatic diester (OAD) groups,\nthe polyoxyethylene terephthalate blocks (TE-POE), the units at the chain ends.\nThe calculation makes it possible to deduce therefrom:\nthe amount by weight of all the (TE) units in the terephthalic polyester composition (TPC),\nthe amount by weight of mono(oxyethyleneoxy) (OEO) residues in the terephthalic polyester composition (TPC).\nFurther information with regard to these analyses is given in the examples."}
-{"text": "This invention relates to respirators in general and, in particular, to a new and useful respirator with an oxygen-releasing chemical cartridge."}
-{"text": "The present invention relates to a thermal ink jet printhead and method of manufacture therefor, and more particularly to an improved thermal ink jet printhead which avoids the effects of standoff between two bonded parts thereof caused by topographic formations developed during fabrication in a thick film insulating layer sandwiched between said two parts.\nIn existing thermal ink jet printing systems, an ink jet printhead expels ink droplets on demand by the selective application of a current pulse to a thermal energy generator, usually a resistor, located in capillary-filled, parallel ink channels a predetermined distance upstream from the channel nozzles or orifices. U.S. Pat. No. Re. 32,572 to Hawkins et al. exemplifies such a thermal ink jet printhead and several fabricating processes therefor. Each printhead is composed of two parts aligned and bonded together. One part is a substantially flat substrate which contains on the surface thereof a linear array of heating elements and addressing elements (heater plate), and the second part is a substrate having at least one recess anisotropically etched therein to serve as an ink supply manifold when the two parts are bonded together (channel plate). A linear array of parallel grooves are also formed in the second part, so that one end of the grooves communicate with the manifold recess and the other ends are open for use as ink droplet expelling nozzles. Many printheads can be made simultaneously by producing a plurality of sets of heating element arrays with their addressing elements on a silicon wafer and by mating a second silicon wafer having a corresponding plurality of sets of channel grooves and associated manifolds therein. After the two wafers are aligned and bonded together, the mated wafers are diced into many separate printheads.\nImprovements to such two part thermal ink jet printheads include U.S. Pat. No. 4,638,337 to Torpey et al. that discloses an improved printhead similar to that of Hawkins et al., but has each of its heating elements located in a recess (termed heater pit). The recess walls containing the heating elements prevent lateral movement of the bubbles through the nozzle and, therefore, the sudden release of vaporized ink to the atmosphere, known as blow-out, which causes ingestion of air and interrupts the printhead operation whenever this event occurs. In this patent, a thick film insulative layer such as polyimide, Riston.RTM. or Vacrel.RTM. is formed on the wafer containing the heating element and patterned to provide the recesses for the heating elements, so that the thick film layer is interposed between the two wafers when they are mated together. U.S. Pat. No. 4,774,530 to Hawkins further refines the two part printhead by disclosing an improvement over the patent to Torpey et al. In this patent, further recesses (termed bypass pits) are patterned in the thick film layer to provide a flow path for the ink from the manifold to the channels by enabling the ink to flow around the closed ends of the channels, thereby eliminating the fabrication steps required to open the groove closed ends to the manifold recess. The heater plates, having the aforementioned improvements of heater pits and bypass pits formed in the thick film insulative layer covering the heater plate surface, are aligned with and bonded to the channel plate, so that each channel groove has a recessed heating element therein and a bypass pit to provide an ink passage from the ink manifold to the channel groove.\nThorough bonding between heater and channel plates is paramount to maintaining the printing efficiency, droplet size consistency, and operational reliability of an ink jet printhead. U.S. Pat. No. 4,678,529 to Drake et al. discloses a method of bonding ink jet printhead components together by spin coating or spraying a relatively thin, uniform layer of adhesive on a flexible substrate and then manually placing the flexible substrate surface with the adhesive layer against the channel wafer surface having the etched sets of channel grooves. A uniform pressure and temperature is applied to ensure adhesive contact with all coplanar surface portions and then the flexible substrate peeled away, leaving a uniformly thin coating on the channel wafer surface to be bonded to the heater wafer. This labor intensive method tends to permit adhesive layer thickness variation between wafer pairs as well as between different regions of the same wafer pair, so that after the wafer pairs are diced into a plurality of identical printheads, the ink flow characteristics varies from printhead to printhead. A more mechanized process to place the adhesive coating on the channel wafer without operator involvement and consequent variation in parameters which introduced thickness variations in the amount of adhesive layer transferred to the channel wafers, especially in the thickness variations from wafer-to-wafer, is described in U.S. patent application Ser. No. 07/888,220, to Narang et al., Filed May 26, 1992.\nAlthough advances have improved the thickness uniformity of the adhesive layer which bonds the ink jet printhead heater and channel plates, insufficient adhesion between bonded heater and channel plates continues to cause a host of problems affecting printhead operation, such as, for example, different drop sizes between adjacent channels, because unwanted topographical protruding formations or lips are formed in the thick film layer during the patterning of the heater pits and bypass pits which prevent adequate contact between the channel wafer surface with the adhesive layer and the thick film layer on the heater wafer. Since increased adhesive layer thickness is not a practical solution because it tends to spread or wick into the channels, the inter-channel gaps between bonded heater and channel plates must be minimized in order to insure consistent printhead firing characteristics. As taught by the above identified U.S. patents, two wafers are bonded together after alignment for subsequent dicing into individual printheads. Each printhead part is formed individually on two separate substrates or wafers, where one contains heating elements and the other ink channels or passageways. The wafer containing the ink channels is silicon, and the channels are formed by an anisotropic etching process. The anisotropic or orientation dependent etching has been shown to be a high yielding process that produces very planar and highly precise channel plates. The other wafer containing the heating elements as well as heater addressing logic is covered by a thick film insulating layer in which heater and bypass pits are formed using photolithography. The thick film layer is preferably polyimide, because it is impervious to water, a major component of the printhead ink. However, one drawback with the polyimide material is its tendency to form unwanted topographical formations, such as raised edges or lips (1-3 microns high) at photoimaged edges. When bonding both heater and channel plates together, a standoff between the two plates is caused by the raised edges, which reduces the adhesiveness of the bond between the two plates and which cause the formation of inter-channel gaps.\nPolyimide topography, such as raised edges or lips, are undesirable by-products resulting from photoimaged and cured heater pits and bypass pits or trenches on heater plates. The raised edges are polyimide topographical features that are formed at the edge of photoimaged areas that do not shrink during curing as would bulk (or centered) areas of the polyimide. Consequently, raised edges critically interfere with the mating and bonding of the heater and channel plates.\nAnother form of polyimide topography is encountered in the form of edge beads or raised areas at the edge of the wafer, when a layer of liquid polyimide is dispensed and spun onto a wafer. When the contact area on the wafer is incapable of spreading further due to the contact angle at the edge of the wafer, centripetal forces push the spinning liquid polyimide towards the outside of the wafer to form an edge bead. The edge bead on a 4 inch diameter wafer, for example, is on the order of 0.5 inch wide radially from the outer edge thereof and has a thickness several micrometers thicker than the rest of the polyimide layer. Such edge beads of polyimide prevent adequate bonding between the wafers. Edge beads can also cause printhead reliability problems, because of additional stress placed on the center area of the channel plate during heater and channel plated bonding may cause cracking. Edge beads, if removed from the edge of the heater wafers, cantilevers the channel plate at its outside edges and can again cause cracks to be formed in the outer peripheral area of the channel wafer. Such cracking in the channel wafer will degrade the reliability of the individual printheads after they have separated from the wafer pair.\nRaised edges and edge beads, however, are not the only topographical formation created from photoimaged polyimide. Other topographical formations, such as wall sags or dips, compound the negative effects of raised edges by adding to the standoff between the bonded heater and channel plates. Wall dips are slumps in the polyimide walls between closely adjacent polyimide photoimaged pits. The polyimide sandwiched between the two wafers can form more than 2 microns of topographical variation, which does not allow the bonding adhesive, approximately 2 microns or less thick, to bridge or fill in the formation of inter-channel gaps. These inter-channel gaps can allow crosstalk between channels when drops are being ejected. As the patent '529 to Drake et al. teaches, care must be taken when applying adhesive in bonding the channel and heater plates so as to insure all surfaces in contact with the ink are free of adhesive, in order that the ink channels are not obstructed during operation.\nOne method of minimizing heater and channel plate standoff of printheads using a modified printhead fabrication sequence is disclosed in U.S. patent application Ser. No. 07/997,473, entitled \"Ink Jet Printhead Having Compensation For Topographical Formations Developed During Fabrication\", assigned to the same assignee as the present invention and filed on Dec. 28, 1992. The printhead enables better bonding of the two plates by compensating for raised lips or edges formed on the outside edge of opposing last pits in an array of pits located in the thick film layer that are created while photofabricating the pits in the insulating layer. The fabrication sequence compensates for the raised edges by including a non-functional straddling channel that nullifies the standoff created by the raised edge and a corresponding additional non-functional pit that positions the raised edge away from the functional channels and nozzles. Although this fabrication technique compensates for polyimide raised edges, it does not attempt to solve the problem of edge bead.\nThere continues to exist, therefore, a need to prevent the standoff between mated heater and channel plates caused by edge beads and without requiring extra non-functional, straddling channels or in drastically altering the fabrication sequence of the heater and channel plates."}
-{"text": "Aspects of this disclosure relate generally to telecommunications, and more particularly to techniques for sounding sequence protection in multi-user multiple-input-multiple-output (MU-MIMO) communications for wireless local area networks (WLANs).\nThe deployment of WLANs in the home, the office, and various public facilities is commonplace today. Such networks typically employ a wireless access point (AP) that connects a number of wireless stations (STAs) in a specific locality (e.g., home, office, public facility, etc.) to another network, such as the Internet or the like. A set of STAs can communicate with each other through a common AP in what is referred to as a basic service set (BSS). Nearby BSSs may have overlapping coverage areas and such BSSs may be referred to as overlapping BSSs or OBSSs. In some scenarios, communications that occur in nearby BSSs can result in collisions and failure in the transmission of information.\nIn dense enterprise deployments of WLANs, such as in stadiums, airports, or other large venues, for example, there may be multiple APs deployed, and the coverage of several of those APs can overlap creating OBSS scenarios. For example, in these dense deployments, multiple STAs can be in the common coverage of multiple BSSs. Moreover, when these dense deployments are unplanned, some of the APs may be automatically configured to work on the same channel, which may cause transmission collisions between OBSSs. The collisions that occur may result in sounding sequence failures and, upon detection of a sounding sequence failure, an AP may terminate a transmission opportunity (TxOP) and would need to contend for the medium again. System throughput can be severely impacted if this happens frequently.\nFor MU-MIMO transmissions, however, the WLAN standards (e.g., IEEE 802.11-based standards) have not defined a specific mechanism that may be used for MU-MIMO sounding sequence protection. One option for MU PLCP Protocol Data Unit (PPDU) protection is to send multiple RTSs one by one to all of the STAs that will be part of the subsequent MU PPDU transmission and expect a CTS from each of those STAs. However, this solution may not be practical for MU-MIMO sounding sequence protection because it incurs a large overhead that may not be needed in non-OBSS scenarios.\nAccordingly, for scenarios that create OBSSs between nearby BSSs and that can result in sounding sequence collisions and failures, it may be desirable to have a mechanism that protects the sounding sequence while providing low overhead."}
-{"text": "1. Field of the Invention\nThe subject invention is directed to gas turbines, and more particularly, to a system for delivering fuel to the combustion chamber of a gas turbine engine by lean direct injection.\n2. Background of the Related Art\nWith increased regulation of pollutants from gas turbine engines, a number of concepts have been developed to reduce engine emissions while improving engine efficiency and overall operability. One such concept is the use of staged combustion. Here, the combustion process is divided into two or more stages or zones, which are generally separated from each other, either radially or axially, but still permitted some measure of interaction. For example, the combustion process may be divided into a pilot combustion stage and a main combustion stage. Each stage is designed to provide a certain range of operability, while maintaining control over the levels of pollutant formation. For low power operation, only the pilot stage is active. For higher power conditions, both the pilot and main stages may be active. In this way, proper fuel-to-air ratios can be controlled for efficient combustion, reduced emissions, and good stability.\nIn addition to staged combustion, providing a thoroughly blended fuel-air mixture prior to combustion, wherein the fuel-to-air ratio is below the stoichiometric level so that combustion occurs at lean conditions, can significantly reduce engine emissions. Lean burning results in lower flame temperatures than would occur during stoichiometric combustion. Since the production of NOx is a strong function of temperature, a reduced flame temperature results in lower levels of NOx. The concept of directly injecting liquid fuel into the combustion chamber of a gas turbine and enabling rapid mixing with air at lean fuel-to-air ratios is called lean direct injection (LDI).\nThe prior art is replete with example of LDI systems. For example, U.S. Pat. No. 6,389,815 Hura et al. discloses a lean direct injection system, which utilizes radially staged combustion within a single injector. The pilot fuel delivery stage includes a pressure swirl atomizer that sprays liquid fuel onto a filming surface. The liquid film is then stripped off into droplets by the action of compressor discharge air. The main fuel delivery system includes a series of discrete atomizers that spray fuel radially outward into a swirling cross-flow of air. The main fuel delivery system is staged radially outboard of the pilot fuel delivery system, and operates in the fuel-lean mode. Radial separation as well as an air jet located radially between the two stages achieves separation of the pilot combustion zone and the main combustion zone.\nU.S. Pat. No. 6,272,840 Crocker et al. discloses a lean direct injection system, which also utilizes radially staged combustion within a single injector. The pilot fuel delivery is either a simplex air-blast type atomizer or a prefilming air-blast type atomizer, and the main fuel delivery system is a prefilming air-blast type atomizer. Separation of the pilot and main combustion zones is achieved by providing an air splitter between the pilot outer air swirler and the main inner air swirler. The air splitter develops a bifurcated recirculation zone that separates the axially aft flow of the pilot injector from the axially aft flow of the main injector. The bifurcated recirculation zone aerodynamically isolates the pilot flame from the main flame, and ensures that the pilot combustion zone remains on-axis with no central recirculation zone. A converging wall of the pilot air cap, which essentially acts as a flame holder to anchor the flame, defines the air splitter. Acting in this manner, the pilot air cap will likely suffer thermal distress (i.e., oxidation, melting), and require some form of thermal management. In this regard, Crocker et al. disclose the use of small cooling holes in the air cap to improve durability.\nEuropean Patent Application EP 1413830 A2 discloses a lean direct injection system, which also utilizes radially staged combustion. In this case, an air splitter with an aft end cone angled radially outward assists in creating a bifurcated recirculation zone. The additional function of the splitter is to prevent the inner main air stream from modulating with combustor pressure fluctuations, thus reducing combustion instability. This air splitter has a larger radial extent than the air splitter disclosed in U.S. Pat. No. 6,272,840 to Crocker et al., and acts as an even larger flame-holder, requiring thermal management to avoid distress.\nWhile the concept of the LDI system is sound, achieving the required levels of performance can be difficult. Lean-burning systems are prone to localized flame extinction and re-ignition. This results in combustion instability that can damage the combustion chamber. Limitations in atomization, vaporization, and fuel-air mixing can result in heterogeneous stoichiometric burning, which yield higher than desired levels of NOx. Also, for these self-contained radially staged LDI systems, control over the level of mixing between the pilot combustion zone and the main combustion zone can be difficult. The negative effects can include reduced margin for lean blowout, and possibly increased levels of smoke.\nAccordingly, there is a continuing need in the art to provide a lean direct injection system which can achieve low levels of combustion instability, enhanced atomization quality, increased fuel-air mixing rates, low pollutant formation, low smoke and improved lean blow-out margin."}
-{"text": "FIG. 1 is a drawing illustrating a semiconductor device in the form of a \u201csystem in a package\u201d (SiP) obtained through a related art semiconductor manufacturing method.\nReferring to FIG. 1, a related art SiP semiconductor device may include interposer 11, first device 13, second device 15, and third device 17.\nFirst device 13, second device 15, or third device 17 may include any one of a central processing unit (CPU), a static random access memory (SRAM), a dynamic random access memory (DRAM), a flash memory, a logic large scale integrated circuit (LSI), a power integrated circuit (IC), a control integrated circuit (IC), an analog large scale integrated circuit (LSI), a microwave monolithic integrated circuit (MM IC), a complementary metal-oxide semiconductor radio-frequency integrated circuit (CMOS RF-IC), a sensor chip, a micro-electro-mechanical systems (MEMS) chip, etc.\nConnection units may be formed between first and second devices 13 and 15, and between second and third devices 15 and 17, and may transfer signals to each device.\nHowever, to extensively use the SiP semiconductor device having the above structure, a problem regarding heat dissipation may need to be solved. Particularly, the problem regarding heat dissipation of a device, such as second device 15 positioned in a mid-layer, may need to be solved to extensively use the SiP semiconductor device."}
-{"text": "Conventional paper shredding machines mostly include a roller blade set constructed of two roller blades that shred or cut paper to be fed into strips or scraps as a result of the opposed rotations of the two roller blades such that information as recorded on the paper is destroyed for confidentiality, and the strips or scraps of paper can be easily compressed to reduce processing space. However, optical discs, regardless of CD-ROM discs or CD-R discs, rather than paper have evolved to be one of the major means for storing information. Once information contained in such optical discs has lost its original value and needs to be destroyed, manually breaking the optical discs not only cannot destroy the information as stored, it also may cause personal injuries. An optimum measure is to shred these discs by means of mechanical operations such as those in conventional paper shredding machines.\nIn the highly economized society as of today, plastic money, such as credit cards, debit cards, ATM cards, or even membership cards issued by enterprises for promotional purposes, and registration cards issued by medical institutions, have made \u201ccards\u201d become an article that can certainly be found in everyone's pocket. When these cards have expired or been replaced with new cards, the most commonly adopted measure is to cut the cards in halves for disposal. However, danger still exists in such a disposing measure because most cards carry the user's signature and the registration cards may also carry personal, medical history, or personal information. It is possible that other individuals with malicious intention may still have access to these halved cards."}
-{"text": "1. Field of the Invention\nThe present invention relates to a fixed-focus imaging lens that forms an optical image of a subject on an imaging device, such as a charge coupled device (CCD) and a complementary metal oxide semiconductor (CMOS), and to an imaging apparatus, such as a digital still camera, a cellular phone with a camera, a mobile information terminal (PDA: Personal Digital Assistance), a smartphone, a tablet terminal, and a mobile game machine, on which the imaging lens is mounted to perform photography.\n2. Description of the Related Art\nAs personal computers have become popular in homes, digital still cameras which are capable of inputting image information about photographed scenes, persons, and the like into the personal computers have spread rapidly. Further, a cellular phone, a smartphone, or a tablet terminal in which a camera module for inputting images is installed has been increasing. Such apparatus having an imaging function uses an imaging device, such as a CCD and a CMOS. Recently, because the imaging device has been miniaturized, there has been also a demand to miniaturize the whole of the imaging apparatus and an imaging lens mounted thereon. Further, since the number of pixels included in the imaging device has also been increasing, there has been a demand to enhance the resolution and performance of the imaging lens. For example, there has been a demand for performance corresponding to high resolution of 5 megapixels or higher, and preferably performance corresponding to high resolution of 8 megapixels or higher.\nTo satisfy such demands, it can be considered that the imaging lens is composed of five or six lenses, which are a relatively large number of lenses. For example, Japanese Patent No. 4858648 (Patent Document 1) and Japanese Patent No. 4792605 (Patent Document 2) propose an imaging lens composed of five lenses. The imaging lens disclosed in Patent Documents 1 and 2 substantially consists of, in order from an object side, five lenses of a first lens that has a positive refractive power, a second lens that has a negative refractive power, a third lens that has a positive refractive power, a fourth lens that has a positive refractive power, and a fifth lens that has a negative refractive power."}
-{"text": "This invention relates to liquid purification and separation, for example to the dewatering of a slurry or the like from sewage treatment plants, where a flocculation compound is mixed with a slurry. The solids of the mixture are filtered from the liquids of the mixture, and a relatively dry sludge cake is discharged from the system.\nIt is desirable to dewater slurries in various manufacturing processes so as to recapture and use the solids and to dispose of the water in an uncontaminated condition. For example, in the disposal of sewage, a sludge slurry is removed from settling tanks of sewage treatment plants which contains both solids and liquid, with a substantial amount of solids suspended in the liquid. One common procedure for treating sludge is to add a flocculation compound as a filter aid to the sludge to free the trapped liquid and filter the solids from the liquid. The flocculants are added to the slurry for increasing the filterability of the slurry and for coagulating the solids, therefore reducing the tendency of the small solids from passing with the liquid through the filter medium. Even after flocculation of the sludge mixture some of the solids are so small and the quantity of liquid so great that the smaller solids tend to move from the mixture with the liquid through the filter medium if the liquid is urged too vigorously from the mixture. This condition requires that the sludge be carefully handled and gently dewatered.\nIn the past, several sludge dewatering systems have been developed which comprise pairs of moving conveyor belts usually fabricated from woven textile material that function as filters. The prior art systems deposit coagulated sludge on one moving belt, and as the sludge moves with the first moving belt the second moving belt is moved down into contact with the sludge so as to squeeze and therefore dewater the sludge. Examples of such prior art systems are disclosed in U.S. Pat. Nos. 3,800,952, 3,601,039, 3,699,881, and 3,896,030.\nIn a sewage treating process it is desirable to process as much sludge as possible in a short period of time, to extract as much liquid as possible from the sludge during the process to produce a dry sludge cake, and to capture a high percentage of the solids from the original slurry. These three objectives are in conflict with each other since more water can be extracted from sludge if the sludge is given more time to drain. Although the dewatering capacity of a dewatering system can be increased by merely speeding up the system, the increase in speed also usually causes a lower percent of solids to be filtered from the liquid and the sludge is not squeezed for a time long enough to remove the desired amount of liquid. When the draining of the sludge is assisted by squeezing the sludge between sandwiched endless conveyor belts, the smaller sludge particles tend to pass through the conveyor belts. Also, if the sandwiched conveyor belts are run through a sinuous path so as to create shear forces in the sludge filter cake, the shifting of the filter cake tends to break the filter cake and allow some of the finer solids to pass with the liquid through the filter belts. On the other hand, after a high percentage of the liquid has been removed from the sludge cake, the sludge cake can be squeezed and higher compression forces can be applied to the sludge cake to remove more liquid substantially without hazard of allowing the small solid particles to pass through the filter belt, since the reduced quantity of water present in the sludge cake is insufficient to wash out the small particles.\nOnce the sludge has been dewatered, it is desirable to remove the sludge from the pores of the filter belts, so that a porous belt is available on the next cycle of the belt for receiving the sludge slurry. The cleaning of the filter belts is usually accomplished by directing a continuous flow of water against the belt, and requires the excessive use of water.\nDewatered sludge can be used for various purposes, such as for a landfill as a fertilizer, or as a fuel. If the dewatered sludge has a low liquid content, it weighs less for a given volume, which means that it can be transported at a lower cost, and the dry product can be better used in subsequent processes since the energy required to extract the residual water would be less and the dryer the sludge, the better the sludge sustains combustion. Also, a dryer sludge makes a better landfill."}
-{"text": "The present invention relates to a hermetic reciprocating compressor including a hermetic casing, a compressing portion having a cylinder and a piston reciprocating inside the cylinder and a motor for driving the piston.\nA hermetic compressor is generally employed in a cooling system such as a refrigerator or an air conditioner, so as to compress a gaseous refrigerant received from an evaporator and supply the compressed refrigerant to a condenser.\nFIGS. 3 and 4 show front and side sectional views of a conventional hermetic reciprocating compressor, respectively. As shown in the drawings, the conventional compressor includes a casing 101 forming a closed internal space, a driving motor 110 installed inside the casing 101 and a compressing portion 120 which is driven by the driving motor 110 to compress a refrigerant. The driving motor 110 includes a stator 112, a rotor 111 rotatably installed inside the stator 112 and a crank shaft 117 fitted into the rotor 111 and rotating therewith while supported within a support member member. The compressing portion 120 includes a cylinder 113, a piston 123 reciprocating inside the cylinder 113 and a cylinder head 127. The piston 123 is connected to an eccentric portion 119 of the crank shaft 117 by a connecting rod 121 and reciprocates inside the cylinder 113 according to a rotational movement of the rotor 111, thereby inhaling and compressing the refrigerant. A suction muffler 141 is installed on the cylinder head 127 to guide the refrigerant to be compressed into an internal space 125 of the cylinder 113. A suction pipe 129 for transferring the refrigerant from an evaporator (not shown) to the compressor is connected to the suction muffler 141 after passing through a wall of the casing 101.\nReferring to FIG. 5, the suction muffler 141 has an internal space for receiving the refrigerant, an inlet 142 to which the suction pipe 129 is connected and an outlet which communicates with a refrigerant inlet 143 provided in the cylinder head 127. A coil spring 145 is installed between the inlet 142 of the suction muffler 141 and the suction pipe 129. One end of the coil spring 145 is fixedly inserted into the inlet 142 of the suction muffler 141 and the other end thereof is fitted outside the leading end of the suction pipe 129, so as to guide the refrigerant passing through the suction pipe 129 into the suction muffler 141.\nThe gaseous refrigerant from the evaporator flows into the suction muffler 141 via the suction pipe 129 and the coil spring 145, and is then supplied to the internal space 125 of the cylinder 113 through the cylinder head 127. On the other hand, the gaseous refrigerant supplied from the evaporator to the cylinder 113 contains liquid oil for lubrication and rust prevention for internal parts of the refrigerant circulation system. Since the refrigerant is vaporized in the evaporator by absorbing heat from the surroundings, whereas the oil maintains the liquid state due to its having a higher evaporation point than that of the refrigerant, the oil in the liquid state flows together with the gaseous refrigerant. The mixture of the liquid oil and the gaseous refrigerant contained in the internal space 125 of the cylinder 113 may damage the piston 123, the inner wall of the cylinder 113 or a valve plate (not shown) of the compressing portion 120, due to a liquid compression phenomenon of the liquid oil. Also, the liquid oil, having a relatively high specific volume, may obstruct the compression of the gaseous refrigerant, thereby decreasing the compression efficiency of the compressor."}
-{"text": "1. Field of the Invention\nThis invention relates to an audio reproducing apparatus and method, audio recording apparatus and method, audio recording and reproducing system, audio data transmission method, information receiving apparatus, and recording medium which are particularly suitable for use in a headphone stereo, and in particular to those configured to store in a portable headphone stereo desired audio data externally transmitted to the portable headphone stereo.\n2. Description of the Related Art\nAmong portable headphone stereo devices with which a user can enjoy reproduced music either during his walk or in an automobile, most widely distributed are those using an analog-recording magnetic tape such as compact cassette. A user of a portable headphone stereo using an analog-recording medium tape typically records in a compact cassette a desired piece of music selected from FM broadcasting programs, CDs (compact discs) or other music sources, and sets the compact cassette in his portable headphone stereo to enjoy audio reproduction either during his walk or in an automobile.\nHowever, as long as compact cassettes or other analog-recording magnetic tapes are used, great improvements in quality of sound are not prospective, and dubbing causes deterioration of signals. Moreover, it takes a long time for a user to record desired pieces of music from CD or other music sources. Additionally, with compact cassettes or other magnetic tapes, the access time is slow, and a user cannot quickly search out, reproduce or repeat his desired music.\nSome portable headphone stereos use CDs. Since CDs are media exclusive for recording, a user of portable headphone CD stereo buys CD recorded with his desired music, and places the CD in his portable headphone CD stereo to enjoy audio reproduction during his walk or in an automobile. CDs are digital recording media, the quality of sound is excellent. The access speed is high, and any desired music can be reproduced quickly. However, since CD headphone stereos are exclusively for reproduction, users cannot make their own music sources compiling their desired pieces of music. Moreover, CD headphone stereos are affected by vibrations, and sound is often skipped over due to external vibrations.\nAlso known are portable headphone stereo players using DAT (digital audio tape), NT (non-tracking tape(trade mark)) or other digital-recording magnetic tapes as their recording media. Deterioration of signals by dubbing hardly occurs in devices using DAT, NT or other digital magnetic tapes. DAT promises audio reproduction of a very high quality of sound. NT permits recording over a long time in a very small cassette. Here again, however, devices using magnetic tapes involve the same problems that the access speed is slow and it takes a long time for repeated reproduction or cue search of a desired piece of music.\nAnother type of portable headphone stereo players uses MD (multi-disk(trade mark)). MDs are media for both recording and reproduction, and users can record their desired pieces of music on MDs from CD or other music sources and can place them in their portable MD headphone stereos to enjoy reproduced music during a walk or in an automobile. MDs are digital-recording media, and promise an excellent quality of sound. The access speed is high, and any desired music can be reproduced quickly. As a shockproof memory is used, devices are less affected by external vibrations.\nAs reviewed above, various kinds of recording media have been proposed for use in portable headphone stereos. However, none of these recording media used in conventional portable headphone stereos meet all requirements from the viewpoints of easy use and quality of sound.\nDevices using analog-recording compact cassettes have a problem in quality of sound. Those using DAT, NT or other digital-recording magnetic tapes have a problem in access speed. CDs are exclusively for reproduction and weak against vibrations. MDs can be used for both recording and reproduction and can be miniaturized but, since the number of titles of MD on sale is not abundant, it is sometimes difficult for users to obtain their desired music. It takes time to dub CD or other music sources.\nAnother problem with MDs is the use of ATRAC (Adaptive Transform Acoustic Coding) compression and expansion ICs or other ICs developed for exclusive use in MDs and the use of a particular architecture, i.e., a particular circuit arrangement as a method for actually mounting the ICs in order to reduce the entire dimension and decrease the cost. Therefore, such devices cannot be used in different ways, namely, for recording and reproducing a recording medium other than MDs, for example.\nMoreover, tastes of users for music are diverged more and more, and the fashion of music largely changes in a short time. It is difficult to exactly catch the fashion of music with conventional music recording media.\nTaking it into account, the present Applicant previously proposed a portable headphone stereo disclosed in Japanese Patent Laid-Open No. hei 06-131371 U.S. Ser. No. 08/131,943, which is configured to write music information in semiconductor memory used as a recording medium to enable reproduction of the music information any time. However, its interior circuit arrangement comprises an exclusive IC and an exclusive architecture, and as in the case of MDs, cannot realize wider use and extensive use of the device."}
-{"text": "1. Field of the Invention\nThe present invention relates to data acquisition systems, and more particularly to a digitally configurable multiplexer/de-multiplexer for routing a number of analog signals to or from selected analog front ends of a digitization block of a data acquisition system.\n2. Description of the Related Art\nAnalog multiplexers/de-multiplexers (MUX/DEMUX) are typically configured as D:1 devices (e.g., 2:1, 4:1, 8:1 or 16:1) with \u201cN\u201d address inputs for selecting from among \u201cD\u201d analog data signals in which D=2N. An analog MUX/DEMUX is typically bidirectional and may be operated as a MUX to select from among multiple inputs to an output or as a DEMUX to route an analog signal to redirect to a selected one of multiple outputs as known to those skilled in the art. The term \u201cmultiplexer\u201d or \u201cMUX\u201d as used herein is intended to connote bidirectional operation in which the MUX may be used for MUX or DEMUX operations. Each configuration has traditionally been implemented on a separate integrated circuit (IC) or chip. If an analog signal needed to be dynamically re-routed to another path, additional multiplexers or switches were required and were placed in series with the signal path. Such additional devices added cost and complexity to the overall data acquisition system. In many cases, different configurations are required in the same system, which forced the designer to implement and purchase different ICs that suited each of the configurations. For example, some systems required the same signal to be routed to any number of different analog front ends that performed various functions (e.g., filtering, amplifying, etc.) depending upon the characteristics of the signal.\nConventional analog multiplexers are relatively inflexible and consume a significant amount of printed circuit board (PCB) area of the data acquisition system."}
-{"text": "The access point (AP) should be able to support a large number of associated stations (STAs) (e.g. over 6000) which may operate on a very strict energy budget. Such devices could be battery powered sensors which transmit and receive data very rarely and stay in the low power operation mode for relatively long periods of time. The basic operation mode is the distributed coordination function (DCF) mode. In order to support a large number of STAs operating in random access mode the access point may utilize specific techniques to restrict contention to the channel to avoid of collisions of simultaneous transmissions. One such technique to reduce contention and collisions of transmission is a method of grouping STAs to multiple groups and assigning certain parameters for each group to indicate the specific group that can access content for the channel. Such grouping information and parameters for the operation can be delivered to the STAs in an association phase or in broadcast information such as beacons.\nThe STAs may operate in low power mode for prolonged periods of time and therefore the grouping related parameters may be not be valid anymore when a STA wakes up and resumes the channel access operation. Broadcast messages such as beacons are transmitted only occasionally and if a STA stays in low power mode for extended periods of time it may lose synchronization (due to the internal clock drift in the AP and in the STA) and is not able to estimate when the next beacon would be transmitted. This may cause the STA to stay awake for relatively long period trying to receive a beacon transmission.\nIn 802.11 protocol the AP buffers data frames if the STA is in the low power state. The AP informs the STAs about the buffered frames by indicating buffered data in the traffic indication map (TIM) which is transmitted in the beacon message. Once the STA awakes it can retrieve the buffered frames by transmitting some buffered uplink (UL) data which implicitly notifies the AP that the STA is awake, or the STA can transmit a power save poll (PS-poll) to indicate to the AP that it is awake and ready to receive data. Alternatively any frame that is classified to be a so called \u2018trigger frame\u2019 indicates that STA is awake.\nTo improve data exchange between AP and STAs, speed frame exchange can be initiated between AP and STA (or between STA and another STA). AP may have data in the DL buffer but is not able to transmit all the data in the same session since it may have some pending management frames in the buffer (such as a beacon transmission) or other control/data frames it may need to deliver."}
-{"text": "Improvements in computer processing power and broadband technology have led to the development of interactive systems for navigating imagery, such as geographic imagery. Some interactive map navigation systems provide a user interface with navigation controls for navigating cities, neighborhoods, and other terrain in two or three dimensions. The navigation controls enable users to tilt, pan, rotate, zoom, and activate terrain and buildings for different perspectives at an area of interest. Exemplary map imaging systems include the Google Earth and Google Maps systems by Google Inc.\nThe imagery used by these interactive systems can be derived by rending geometric objects and texture objects to generate a two-dimensional or three-dimensional graphical representation of an area of interest. The geometric objects can define three-dimensional surfaces of objects and terrain depicted in the imagery. The texture objects can be mapped to the three-dimensional surfaces to add detail, surface texture, color and other features to the depicted objects and terrain.\nInteractive systems for displaying imagery, such as geographic imagery, often render geometric and texture objects with low level of detail (\u201cLOD\u201d) (e.g. low resolution) at camera views far from an area of interest to improve performance. As the user zooms in closer to the area of interest, geometric and texture objects with higher LOD (e.g. higher resolution) are rendered to depict objects in greater detail. Switching between two levels of detail can cause a distracting pop if the levels of detail are sufficiently different and/or the transition happens too close to the camera view. The pop can be a visual anomaly which disrupts the user's sense of immersion in the imagery. Moreover, if the geometric data and the texture data are fetched from a server over a network, the latency between the request for higher resolution objects and the receipt of the data by the client can be long. In this case, a user can zoom in past the natural resolution of a rendered scene, resulting in a blurry image. When the higher resolution imagery finally arrives from the server and is rendered, the higher resolution imagery will \u201cpop\u201d into view because its resolution is markedly different from the imagery in the previously displayed blurry image."}
-{"text": "The invention relates to a retractable blade utility knife handle and more particularly to a retractable blade utility knife handle having a pair of releasably locking, pivotable handle halves.\nIn retractable-blade utility knives, the blade is slidably movable from a sheathed position to an extended unsheathed position wherein the knife blade projects through an opening in the knife handle to present a cutting edge. Such a utility knife is disclosed in Weimann U.S. Pat. No. 4,663,845, entitled \"Utility Knife\" wherein the knife incorporates a blade carrier which mounts and supports the blade within the interior of the knife handle for selective longitudinal movement therein. The blade carrier includes a thumb actuated button to release the blade carrier from one of several latching positions for slidably shifting the blade to another longitudinal position. The halves of the knife handle are connected by a central pivot. The handle halves are releasably locked so that they can be pivoted laterally relative to one another when the blade carrier is in any longitudinally shifted position, thereby allowing the blade to be replaced and providing access to a blade storage compartment.\nA typical shortcoming of retractable-blade utility knives is the lack of stability of the blade when the knife is forcibly twisted sideways during use. For example, under conditions in which a substantial force is applied to the side of the blade proximate the lower cutting edge, it is possible for the lower edge of the blade to slip off of the blade carrier, thereby freeing the blade to fall out of the handle. Another drawback of conventional utility knife handles is the possibility that the handles can be opened inadvertently. Yet another shortcoming of conventional utility knives is that the blade carrier and blade tend to rattle, as they are not fixed to the handle. A further drawback of conventional retractable blade knives is that if the blade carrier is shifted longitudinally when no blade is installed in the blade carrier, the corners of the blade carrier, on occasion, will become stuck at various locations in the handle, causing the blade carrier to shift or slide irregularly."}
-{"text": "This invention relates to novel carboxycycloalkylamino derivatives. The carboxycycloalkylamino derivatives of the present invention are modulators of the sphingosine-1-phosphate (S1P) receptors and have a number of therapeutic applications, particularly in the treatment of hyperproliferative and autoimmune diseases, in mammals, especially humans, and to pharmaceutical compositions containing such compounds.\nThe S1P receptors 1-5 constitute a family of seven-transmembrane G-protein coupled receptors. These receptors, referred to as S1P1 to S1P5, are activated via binding by sphingosine-1-phosphate, which is produced by the sphingosine kinase phosphorylation of sphingosine. S1P receptors are cell surface receptors involved in a variety of cellular processes, including cell proliferation and differentiation, cell survival, and cell migration. S1P is found in plasma and a variety of other tissues and exerts autocrine and paracrine effects.\nRecent studies indicate that S1P binds to the S1P1 receptor to promote tumor angiogenesis by supporting the migration, proliferation and survival of endothelial cells (ECs) as they form new vessels within tumors (tumor angiogenesis) (Lee et al., Cell. 99:301-312 (1999) Paik et al., J. Biol. Chem. 276:11830-11837 (2001)). Because S1P is required for optimal activity of multiple proangiogenic factors, modulating S1P1 activation may affect angiogenesis, proliferation, and interfere with tumor neovascularization, vessel maintenance and vascular permeability.\nOther diseases or conditions that may be treated with the compounds of the present invention include organ transplant rejection and inflammatory diseases, which are believed to proceed via modulating the S1P receptors.\nThus, the identification of compounds which modulate the activity of the S1P1 receptor to regulate and modulate abnormal or inappropriate cell proliferation, differentiation, or metabolism is therefore desirable."}
-{"text": "The present invention relates to memory cards, and relates more particularly to a contact terminal for memory card plug connectors which prevents contact errors.\nVarious memory cards have been developed for connection to a card jack on a computer system for storing data. As illustrated in FIG. 1, the card jack, referenced by 1, comprises an insertion slot 11 at the front side, and a set of contact pins 12 for connection to the circuit board of the computer system. The memory card, referenced by 3, has a plug connector 4 for insertion into the insertion slot 11, and a set of contact terminals 41 mounted in the plug connector 4 for connection to the contact pins 12 respectively. As illustrated in FIGS. 2 and 3, each of the contact terminals 41 comprises a curved upper clamping section 411 and a curved lower clamping section 412 for contacting one contact pin 12. Because the contact pin 12 is disposed in contact with the upper clamping section 411 and the lower clamping section 412 at a respective contact point, the contact pin 12 tends to displace, causing a contact error. FIGS. 4 and 5 shows another structure of contact terminal for this purpose. This contact terminal, referenced by 42, comprises two substantially C-shaped clamping portions 420 at two opposite ends, and three longitudinal connecting portions 421, 422, and 423 connected between the C-shaped clamping portions 420. However, this structure of contact terminal is still not satisfactory in function. During the installation of the memory card, much effort must be given to the plug connector of the memory card against the card jack so that the contact terminals 42 of the plug connector of the memory card can be respectively forced into engagement with the contact pins 12 of the card jack. However, when the contact terminals 42 are forced into engagement with the contact pins 12 of the card jack, the contact pins 12 tend to be deformed or broken. Another drawback of this structure of contact terminal is that frequently mounting and dismounting the contact terminal will cause the longitudinal connecting portions to lose their elastic resilient material property. If the longitudinal connecting portions lose their elastic resilient material property, a positive contact between the contact terminal and the respective contact pin cannot be obtained."}
-{"text": "The present invention concerns a device for indexable coding for an electrical connector of the type comprising a correction component having a rear part comprising a right prism of which at least one of the lateral faces constitutes an angular positioning reference surface for indexing as well as an elongated front part extending in a longitudinal direction and having a correction profile.\nIndexable coding devices of this type are currently used in the connection field, in particular for connectors designed for aviation.\nIt is desirable to increase the number of correction positions, which is currently feasible only at the price of a certain complication, thus resulting in a relatively high cost."}
-{"text": "Currently, the number of various crimes such as robbery of articles placed inside cars and illegal entry of another's residence has been increasing, and the consciousness of people for crime prevention has been growing. To meet the demand generated with the increase in consciousness for crime prevention, a variety of crime-prevention goods have been developed.\nFor example, a security mechanism control system which locks and unlocks doors of a vehicle has been developed as an example of a security mechanism control system for a vehicle using radio waves. For example, according to security mechanism control systems disclosed in Patent Reference No. 1 (JP-A-2003-27790, publication date: Jan. 29, 2003), Patent Reference No. 2 (JP-A-2003-27791, publication date: Jan. 29, 2003), and Patent Reference No. 3 (JP-A-2003-301638, publication date: Oct. 24, 2003) or other references, a user can automatically lock or unlock a vehicle only by carrying a portable device and easily check that the vehicle has been locked. In these security mechanism control systems, the portable device transmits signals and a main device receives the signals. Thus, the user carrying the portable device can automatically unlock the vehicle only by approaching the vehicle and lock the vehicle by moving away from the vehicle.\nIn such a conventional type of security mechanism control system which automatically locks and unlocks doors of a vehicle when the user carrying the portable device gets into and out of the vehicle, a plurality of portable devices are prepared and separately used in some cases when the security target such as a vehicle is used by a plurality of users.\nWhen the plural portable devices are used, however, the possibility that any of the portable devices is left inside the vehicle increases. If any of the plural portable devices is left inside the vehicle, the following problem occurs according to the conventional system.\nSince the left portable device is located inside the vehicle, the main device determines that the user as the owner of the portable device is still staying inside the vehicle and does not automatically set the security condition for the vehicle in some cases. When this occurs, the user does not notice that the portable device has been left inside the vehicle and leaves the vehicle in the non-security condition.\nIn another case, each of the plural users has his/her portable device and the security mechanism control system automatically locks the vehicle when any of the plural users gets off the vehicle. In still another case, the user locks the vehicle by directly giving a locking command to the vehicle.\nIn either of these cases, the vehicle is automatically unlocked by the approach of a person other than the users when any of the portable devices is left inside the vehicle. This is because the requirements that a person is approaching the vehicle and that the portable device is positioned near the main device are satisfied, and therefore the locking condition set for the vehicle is automatically cancelled.\nMoreover, when the plural portable devices are used in the system, such a case can happen where a user carrying the portable device moves away from the main device but another user stays near the main body.\nIn this case, the user leaving the vehicle considers that the vehicle has been automatically locked, but in reality the vehicle is not locked since the other user stays near the main device. Thus, the operation of the system depends on the behaviors of the users. There is a possibility that the requirement for unlocking is satisfied while the requirement for locking is satisfied. The above case is a result of the operation selected by the main body which has determined that the vehicle should be unlocked.\nTherefore, when the plural portable devices are used, the main device does not always set the condition as the user desires.\nAccording to the structures shown in Patent Reference Nos. 1 through 3, when the portable device left near the main body or carried by the user is positioned near the main body, the problems mentioned above cannot be solved. This is because the system does not have means for determining whether the portable device near the main body is the device left behind or the portable device carried by the user, nor means for appropriately notifying the user of the result of the determination.\nIn the structures shown in Patent Reference Nos. 1 through 3, therefore, in case of the plural portable devices are used, the vehicle cannot be locked or the locked condition is cancelled against the user's intension while the user is not noticing this situation."}
-{"text": "Field of the Invention\nThe present application relates generally to electrical machines including, but not limited to electrical motors and electrical generators and more particularly, but not by way of limitation to electrical machines that are co-axially arranged with a magnetic gear.\nHistory of the Related Art\nElectrical machines have been utilized for a variety of purposes. Typically, the cost of an electric machine is proportional to the torque with which it must interact. For applications requiring a large torque at a low speed there are two generally accepted options. First, a high speed, low torque machine could be utilized in conjunction with a mechanical gear. Second, a larger low speed, high torque direct-drive machine could be utilized. Mechanical gears are generally viewed as unattractive options as they require considerable maintenance, produce acoustic noise, and have a shorter useful life.\nMagnetic gears have attracted considerable attention as a possible replacement for traditional, mechanical gears. Unlike mechanical gears, which rely on the physical interaction between teeth, magnetic gears create a gearing action through the modulated interaction between magnetic flux generated by two means of inducing magnetic flux with different pole pair counts. Magnetic gears exhibit generally contactless operation and, as such, facilitate reduced maintenance, improved reliability, decreased acoustic noise, and physical isolation between input and output shafts."}
-{"text": "The invention relates to a lubricator for underbalanced drilling.\nThere are two techniques that typically are used to drill a borehole in a formation: an overbalanced drilling technique and underbalanced drilling technique. In overbalanced drilling, fluid in an annulus of a well is used to exert a pressure that is greater than the formation pressure. Thus, the pressure that is exerted by the annulus fluid keeps formation fluids from exiting the well. A drawback to this technique is that mud particles typically are added to the annulus fluid to increase its weight (and thus, increase its downhole pressure), and these mud particles tend to clog up openings in the formation. Thus, the formation may be damaged by overbalanced drilling, and after drilling, cleanup of the well may be needed before production begins. The well may also need to be tested after overbalanced drilling to check for formation damage.\nUnlike overbalanced drilling, underbalanced drilling typically does not damage the formation damage and typically maximizes reservoir inflow. In underbalanced drilling, heavy annulus fluid is not used to suppress the formation pressure. Instead, a blowout preventer, or snubbing unit, is used to seal off the drill string at the surface of the well. However, this arrangement may also present difficulties. For example, when drilling at shallow depths or retrieving the drill string, the upward force from the formation pressure may exceed the weight of the drill string and thus, may force the drill string out of the borehole. As a result, retrieving the drill string may consume a considerable amount of time and present a significant danger.\nThus, there is a continuing need for an arrangement to address one or more of the problems that are stated above.\nIn an embodiment of the invention, a system usable with a subterranean well includes a tubing and a lubricator. The tubing is adapted to receive a drill string in a passageway of the tubing, and the lubricator is located downhole and is connected to the tubing. The lubricator is adapted to be remotely operable from a surface of the well to control fluid communication between the passageway located above the lubricator and a formation located beneath the lubricator.\nIn another embodiment of the invention, an apparatus that is usable with a downhole tool that has a passageway includes a sleeve and a controller. The controller selectively moves the sleeve into the passageway to protect a portion of the downhole tool from a downhole fluid.\nAdvantages and other features of the invention will become apparent from the following description, drawing and claims."}
-{"text": "It is generally known to those skilled in the art of golf-club making that a player's performance level may be enhanced by custom fitting a set of golf clubs to the player. These custom-fit golf clubs are selected, using anyone of a plurality of fitting methods, to improve the player's ability to consistently execute well-struck shots. Most fitting methods require a player to perform one or more golf swings using a single test club, e.g., a 6 iron. The collected data is then extrapolated to the rest of the clubs with the set. However, a player who is fitted into a custom set using a single test club may later discover a non-uniform shot-distance progression between the clubs within the set. Accordingly, the player may have difficulty selecting the appropriate golf clubs from the set for various shots during a golf round."}
-{"text": "This application was originally filed as international application PCT/NL99/00498 on Aug. 4, 1999.\nThe present invention relates to a closing valve for a container or a pipe system, comprising a closing jacket connectable to an opening in the container or forming part of the container, which closing jacket is provided on the inside with a narrowed and a widened portion, and a valve part movable in this closing jacket with a closing element and a clamping element rigidly connected therewith and extending in the closing jacket, which clamping element is provided on at least the outside with a thickening which, when the passage through the closing jacket is sealed by the valve part, is in contact with the widened portion in the closing jacket, while, furthermore, a head part movable back and forth in the closing jacket, optionally to be put into circulation apart, is provided, by means of which the valve part can be moved with respect to the closing jacket such that the passage therethrough can be released and/or closed.\nSuch a container may be a cardboard, metal or plastic packing as well as a compressible plastic bag, a glass, bottle and the like. The container is suitable for gaseous products as well as all kinds of liquids, optionally mixed with gas, both under pressure and under vacuum, highly viscous substances, such as puree, and further granulates, granular material, etc. The closing valve may also be included in a pipe system.\nA closing valve as described in the opening paragraph is known from WO94/29215. In said document a thickening at the outside of the clamping element has the function to lock a head part with a tube relative to a valve part in the form of a closing plug. Because said thickening only temporarily rests against the edge of the thickening, for a short time a counterforce, necessary to lock the head part with the tube relative to the valve part, is exerted. By pushing the close jacket in the first instance said thickening meets the edge of the narrowed part of the closing jacket. At that moment the head part with the tube is still free removable (FIG. 2). When the head part with the tube is pushed further, the head part is locked in the xe2x80x98valvexe2x80x99 part because the edge of the valve part engages a groove between the head part and the tube (FIG. 3). When thereafter the head part with the tube is pushed further again, said thickening being pressed along the edge of the narrowed part of the closing jacket, the container is fully opened as then the outlet opening in the tube is located in the container (FIG. 4). In the latter position there is no locking between the valve part with the head part on the one hand and the closing jacket at the other hand, so that the position of these parts relative to each other is not determined any longer and not adjustable at all. As the valve element is in the form of a closing plug, the outlet opening is brought in the tube, with the risk that when the tube is wrenched off the valve part, the valve part will fall into the container, which then cannot be closed again.\nThe object of the invention is to provide for such a container an efficiently and inexpensively producible closing valve, via which the container can be easily filled and closed and the substances contained therein can be easily poured out or pressed or sucked out of the container, and by means of which closing valve the above problems are obviated.\nTo this end, according to the invention the closing valve is characterized in that, the thickening on the outside of the clamping element is located near the upper part of the clamping element, such that, when the passage through the closing jacket is released, the thickening on the outer side of the clamping element is brought into the narrowed portion of the closing jacket, as a result of which the clamping element reaches a position in which it is engaged by the head part.\nThe closing vale is therefore opened by pushing down the valve part by means of the head part, while by pulling up the head part, which as then engaged the clamping element, the valve part is taken along and the closing valve is gradually closed. The head part can therefore also be designed as a tap and brought into several positions and optionally fixed in these positions with respect to the closing jacket, so that gradually opening and closing the passage by the closing jacket becomes possible.\nAlthough it is sufficient for specific uses to only open the container by means of the head part, in which connection it is not important that the valve part lands in the container, it is important for many other uses that the head part is provided with locking means for preventing the valve part from being pushed completely out of the closing jacket.\nAs appears from the foregoing, the closing valve, in its basic structure, is composed of only three easily producible parts: a closing jacket, a valve part and a head part; moreover, an inexpensive final product can be obtained by, for instance, injection molding these parts from some hard plastic. Of course, it remains possible that the closing jacket forms part of the container and/or that the head part forms part of a draw-off tap.\nFor many uses the head part will be provided with a throughbore which, at the lower end, is in open communication with the inside of the closing jacket, so that, when the closing valve is opened by pushing in the head part, an open communication with the atmosphere results, to which optionally other filling, outlet or pouring means or the like can be connected. In a favorable embodiment the head part is provided with a widened upper part which, when the valve part is pushed out of the closing jacket at least partially, seals the upper side of the closing element.\nIn addition to the three above basic parts, further parts may be provided, such as a sealing element to be arranged between the closing jacket and the head part, in particular an O-ring, to enable a fluid-tight closure. Further sealing elements, for instance around the peripheral edge of the closing element, may be provided; these sealing elements, however, may also be formed by, for instance, plastic edges forming part of the relevant components.\nAlthough the clamping element may be cylindrical, while an annular thickening may be provided on the outside as well as on the inside, it is difficult to push in such a clamping element on all sides; anyhow, the clamping element then needs to be rather easily deformable. It is better when the clamping element is finger-shaped, with each of the finger-shaped parts on the inside and outside being provided with a thickening. It is indeed easier to push in the individual fingers. The thickening on the inside then serves to facilitate engagement by the head part.\nWhen the container is formed by an easily compressible packing, such as, for instance, a plastic bag to which the closing valve as hitherto defined is connected, and furthermore, this easily compressible packing is arranged in an outer packing which is not compressible or is relatively difficult to compress, the closing valve needs to be provided with additional means. To this end, according to the invention a collar element connectable to an outer envelope or forming part thereof is provided, which collar element fits around the closing jacket, while between the outer circumferential surface of the closing jacket and the inner circumferential surface of the collar element gas supply channels, in particular air supply channels, are provided, which are opened when releasing the passage in the closing jacket. This has the advantage that gas can flow into the space between the inner packing, this is the easily compressible packing, and the outer packing, so that, when a pump is connected to the head part, the inner packing can be easily emptied by sucking, while the inner packing is compressed. Also, gas can be forcibly brought into the above space to enable the contents to be pressed out of the inner packing. This is important when the head part is formed by a filling head or a tapping head of, for instance, a tapping plant for beverages."}
-{"text": "(a) Field\nExemplary embodiments of the invention relate to a quantum rod sheet, a backlight unit including the quantum rod sheet, a display device including the backlight unit and a manufacturing method thereof.\n(b) Description of the Related Art\nVarious types of flat panel display, such as a liquid crystal display (\u201cLCD\u201d), a plasma display panel (\u201cPDP\u201d), an organic light emitting display (\u201cOLED\u201d), an electrowetting display (\u201cEWD\u201d), an electrophoretic display (\u201cEPD\u201d), an embedded micro-cavity display (\u201cEMD\u201d), and a nano-crystal display (\u201cNCD)\u201d, have been developed from a cathode ray tube (\u201cCRT\u201d) type using a cathode ray tube.\nAmong the various types of flat panel display, the LCD, which is one of the most widely used types, has characteristics, such as miniaturization, weight lightening, and low power consumption. In general, the LCD is a device, in which an electric field is generated in a liquid crystal material disposed between an upper substrate and a lower substrate by applying different potentials to a pixel electrode and a common electrode thereof such that an arrangement of liquid crystal molecules is changed, and transmittance of light is thereby controlled to display an image.\nIn the LCD, a liquid crystal panel does not emit light therefrom, and the liquid crystal panel thereby includes a backlight unit for providing light to the liquid crystal panel. Among the various types of flat panel display, other light receiving type display devices may include the backlight unit.\nThe backlight unit typically includes a light source and a plurality of optical sheets to improve luminance of light provided from the light source and to distribute the light substantially uniformly over an entire region thereof. Among the optical sheets, an optical sheet that may be used to improve the luminance (e.g., a luminance enhancement film) is typically substantially expensive, which may lead to an increase in manufacturing costs of the display device.\nIn a backlight unit configured to generate white light, the backlight unit typically includes a light emitting diode (\u201cLED\u201d) that emits light having a specific wavelength and a phosphor that changes the light having the specific wavelength emitted from the LED into the white light to provide the white light."}
-{"text": "Occasions arise where a well has to be plugged and abandoned. Typically, such wells will have a cemented casing and another tubular string inside the casing. The plug and abandon procedure involves severing the inner tubular followed by axially milling the severed tubular for a predetermined axial distance followed by setting a plug and depositing cement on top of the plug. The tubular to be severed has a nominal outer dimension between upsets at opposed ends where the joints occur. The axial milling has to cut the nominal dimension as well as the upset outer dimension. In either case it is not desirable to cut into the surrounding tubular as the integrity of the surrounding tubular is important to the functioning of the plug created as part of the abandonment process. On some occasions the tubular to be severed and axially milled is not disposed concentrically in the surrounding casing and sometimes the inner tubular lies on the casing. Axially milling the severed tubular involves using blades that extend a predetermined amount that is short of reaching the surrounding tubular.\nThe way this process is done today takes multiple trips. The first trip is with a mill with extending blades whose pointed ends radially penetrate the wall of a tubular until that wall is breached. These blades are limited in their radial extent to avoid damage to the surrounding tubular. The tool is then tripped out for a blade replacement to dress the tool with blades that will axially mill the nominal outside diameter of the tubular in a location where the severing took place, which is generally between upsets that have an even larger outer dimension. The axial milling continues between the upsets until an upset is reached. At that point there is another trip out of the hole to redress the tool with longer axial milling bladed that can reach further to the exterior of the larger upset dimension and yet not far enough to gouge the surrounding tubular. Should further axial milling be needed after an upset then another trip to swap back to the shorter axial milling blades is needed.\nIn one embodiment of the invention all three blade types are provided on a tool housing run into the tubular to be severed and then axially milled. The severing blades are initially extended for the radial cutting through the wall after which those blades are no longer needed. With the aid of specially configured darts that land in a profile either the shorter or the longer axial milling blades are extended. The uniquely configured darts are blown out into a catcher to allow selection of the longer or the shorter blades alternatively as many times as needed depending on the length of the axial milling distance and how many upsets need milling in that desired interval. No matter if the shorter or longer axial milling blades are extended, there is an open circulation path for the extended blades to cool them and for cuttings removal to a collection device or to the surface. While severing and milling is the preferred application, other uses are envisioned where there is a need to sequentially operate a variety of tools in a single trip particularly when two of the tools need to be alternatively operated multiple times. While the preferred actuation system is described in a hydraulic operating environment other ways to achieve the desired tool operating sequence are also contemplated with the objective being an economical and hence reliable design that can perform the required actions in a single trip and preferably without well intervention.\nU.S. Pat. No. 8,955,597 shows a sequential use of a mill to make a radial cut with blades 76 followed by extension of a reamer 84 to enlarge the opening where the casing was earlier milled away. This application is limited to a single string in the hole and the sequence is executed but a single time. U.S. Pat. No. 5,765,640 describes a multipurpose tool where one or more pistons can be used to operate multiple tools at the same or different times but shows or describes no structure as to how that can be done one time, not to mention multiple times. U.S. Pat. No. 8,141,627 shows multiple rows of cleaning blades that extend at the same time and are spring retracted.\nThose skilled in the art will have a better understanding of the preferred embodiment of the invention from the description below and the associated drawings while recognizing that the full scope of the invention is to be determined from the appended claims."}
-{"text": "This disclosure relates to refrigeration systems for perishable items."}
-{"text": "The present disclosure relates generally computer vision systems and relates more particularly to sensors for measuring the distance to an object or point in space.\nUnmanned vehicles, such as robotic vehicles and drones, typically rely on computer vision systems for obstacle detection and navigation in the surrounding environment. These computer vision systems, in turn, typically rely on various sensors that acquire visual data from the surrounding environment, which the computer vision systems process in order to gather information about the surrounding environment. For instance, data acquired via one or more imaging sensors may be used to determine the distance from the vehicle to a particular object or point in the surrounding environment."}
-{"text": "Sebaceous glands secrete sebum which contain lipids which collect on the skin and in the ear canal. The rate of secretion depends on several factors. Several skin disorders, including acne and seborrhea, are associated with inappropriate sebum production.\nThe removal of excess lipids from the skin can usually be accomplished by bathing using detergents and cleansers. Many preparations have been developed to assist in removal of excess lipids. U.S. Pat. No. 5,026,551 discloses compositions comprising carbon dioxide gas and emulsifiers with oils in bath preparations. In one embodiment the carbon dioxide gas was carried on cyclodextrin. U.S. Pat. No. 4,970,072 discloses use of whey products in bath preparations.\nCyclodextrins have been used as delivery agents for water-insoluble drugs for topical, oral and parenteral delivery. They have also been used to deliver cosmetic preparations to the skin. Several compositions utilize cyclodextrin inclusion products. European patent application 0 366 154 (1990) discloses several cyclodextrin inclusion products for use in cosmetic compositions. U.S. Pat. No. 4,678,598 discloses and claims a cyclodextrin-containing shampoo containing menthol and camphor. The cyclodextrin is provided to depress odor. U.S. Pat. No. 4,267,166 discloses use of cyclodextrin to treat foul breath. U.S. Pat. No. 4,891,361 discloses use of a kojic/cyclodextrin inclusion complexes to prevent elastosis in an animal test by preventing formation of melanin. A publication of Wacker Chemicals discloses that empty cyclodextrins in powder or creams may contain excreted matter of the skin or reaction substances produced on the skin. The statement by Wacker is under a subheading, \"masking of disagreeable smells\" and contains four other statements, all of which pertain to odorous substances. The sebaceous skin lipids are not odorous. Therefore, is appears that the Wacker publication was referring to the containment of unpleasant odors that are produced when bacteria act upon lipids and other skin secretions. There is no evidence therein the materials were useful as cleansing agents for delipidation of the skin."}
-{"text": "Japanese Patent Laying-open No. 2003-304604 discloses a device driving a motor mounted as a source driving electric vehicles, hybrid vehicles, fuel cell vehicles and other various types of vehicles.\nIn this device, if a temperature detection means detects that the motor has a temperature equal to or higher than a temperature limit, a motor output control means limits the output of the motor. Herein a means for detecting a rate of change in temperature detects the rate of change in temperature of the motor and modifies a setting of the temperature limit in accordance with the rate of change detected.\nMore specifically, if the motor's rate of change in temperature is equal to or larger than a predetermined rate of change, the means for detecting a rate of change in temperature determines that the motor has a large increase in temperature, and accordingly the means sets the current temperature limit to be a first temperature limit, and when the motor attains the first temperature limit or higher the motor output control means limits the output of the motor.\nIn contrast, if the motor's rate of change in temperature is smaller than the predetermined rate of change, the means for detecting a rate of change in temperature determines that the motor has a small increase in temperature, and accordingly the means sets the current temperature limit to be a second temperature limit higher than the first temperature limit, and when the motor attains the second temperature limit or higher the motor output control means limits the output of the motor.\nThis device allows a vehicle to run an increased distance without limiting the output of the motor, can protect the motor from high temperature, and also allows the motor to fully exhibit its performance.\nHowever, the motor drive device disclosed in Japanese Patent Laying-open No. 2003-304604 only limits the output of the motor in accordance with the rate of change in the temperature of the motor in view of protecting the motor from high temperature. It does not give consideration to the state of the vehicle.\nFor example, when a vehicle is in a state in which while it is running up hill it may fall downhill, such state should first and most of all be avoided. The above described motor drive device, despite the vehicle's such state, may prioritize protecting the motor from high temperature and stringently limit the motor's output. Thus with the above described motor drive device there is a possibility that while the vehicle is running uphill it may fall downhill."}
-{"text": "The present invention is generally related to medical devices, kits, and methods. More specifically, the present invention provides a system for crossing stenosis, partial occlusions, or total occlusions in a patient's body.\nCardiovascular disease frequently arises from the accumulation of atheromatous material on the inner walls of vascular lumens, particularly arterial lumens of the coronary and other vasculature, resulting in a condition known as atherosclerosis. Atheromatous and other vascular deposits restrict blood flow and can cause ischemia which, in acute cases, can result in myocardial infarction or a heart attack. Atheromatous deposits can have widely varying properties, with some deposits being relatively soft and others being fibrous and/or calcified. In the latter case, the deposits are frequently referred to as plaque. Atherosclerosis occurs naturally as a result of aging, but may also be aggravated by factors such as diet, hypertension, heredity, vascular injury, and the like.\nAtherosclerosis can be treated in a variety of ways, including drugs, bypass surgery, and a variety of catheter-based approaches which rely on intravascular widening or removal of the atheromatous or other material occluding the blood vessel. Particular catheter-based interventions include angioplasty, atherectomy, laser ablation, stenting, and the like. For the most part, the catheters used for these interventions must be introduced over a guidewire, and the guidewire must be placed across the lesion prior to catheter placement. Initial guidewire placement, however, can be difficult or impossible in tortuous regions of the vasculature. Moreover, it can be equally difficult if the lesion is total or near total, i.e. the lesion occludes the blood vessel lumen to such an extent that the guidewire cannot be advanced across.\nTo overcome this difficulty, forward-cutting atherectomy catheters have been proposed. Such catheters usually can have a forwardly disposed blade (U.S. Pat. No. 4,926,858) or rotating burr (U.S. Pat. No. 4,445,509). While effective in some cases, these catheter systems, even with a separate guidewire, have great difficulty in traversing through the small and tortuous body lumens of the patients and reaching the target site.\nFor these reasons, it is desired to provide devices, kits, and methods which can access small, tortuous regions of the vasculature and which can remove atheromatous, thrombotic, and other occluding materials from within blood vessels. In particular, it is desired to provide atherectomy systems which can pass through partial occlusions, total occlusions, stenosis, and be able to macerate blood clots or thrombotic material. It is further desirable that the atherectomy system have the ability to infuse and aspirate fluids before, during, or after crossing the lesion. At least some of these objectives will be met by the devices and methods of the present invention described hereinafter and in the claims."}
-{"text": "The present invention relates to a filter medium for filtering impurities out of fluids and a method of manufacturing the filter medium.\nA filter medium for filtering impurities having relatively large diameters is generally formed by a kind of sieve. However, in a fluid there are contained not only solid grains having large diameters, but also particles having very small diameters and fluids of different kinds. For instance, in an engine oil for use in motor cars there are included various kinds of particles having different sizes such as iron powders, carbon particles, pigments and liquids such as water. Particularly, fine particles and liquids could not be easily removed by the conventional filter medium. Therefore, absorbing agents such as active carbon and activated white clay(kaoline) are provided in the filter medium.\nHeretofore, there has been proposed a charcoal filter in which absorbing agents are mixed in a filtering member or inserted between sheet-like filtering members such as filtering paper and nonwoven fabric. In these known filter medium, it is quite difficult to perform the desired filtering action due to the fact that the absorbing agents might move in the filtering members or flow out of the filtering members.\nIt has been further known that the hot active carbon particles are spread on an nonwoven fabric made of polypropylene resin and are fused onto the fabric. However, in such a filter medium, the absorbing agents, i.e. carbon particles are not firmly secured to the filtering member i.e. polypropylene unwover fabric and further the filterability of the fabric is largely lost at positions where the carbon particles are fused, so that the filtering efficiency is decreased. It has been also known to secure the absorbing agents to the filtering member by means of a bonding agent. In such a known filter medium, the filtering efficiency is also decreased. Further, the bonding agent might affect the absorbing function of the absorbing agents.\nIn an oil cleaner, fine magnetic particals such as iron particles produced by the friction of machine parts should be removed. Such magnetic particles may be absorbed by means of a filter medium including permanent magnet particles. However, it is almost impossible to spread uniformly the magnet particles or powder on the filtering member, because the magnet particles are attracted to each other.\nFurther, in case of filtering impurities out of an oil, it is necessary to remove water contained in the oil. However, in the known filter medium since the filtering member is clogged due to the water absorption by means of OH-radical, the filterability is decreased to a large extent."}
-{"text": "1. Field of the Invention\nThe invention relates to a furniture assembly, more particularly to a furniture assembly, such as a chair, that is easy to assemble and that provides good support.\n2. Description of the Related Art\nGenerally, a conventional inflatable mattress or sofa bed relies only on air introduced into an inner portion thereof to maintain its shape and to support the weight of a user. When the air introduced into the conventional inflatable mattress or sofa bed is insufficient, the conventional inflatable mattress or sofa bed is uncomfortable and unsafe to use."}
-{"text": "It has become common for users of host computers connected to the World Wide Web (the \"Web\") to employ Web browsers and search engines to locate Web pages having specific content of interest to users. A search engine, such as Digital Equipment Corporation's AltaVista search engine, indexes hundreds of millions of Web pages maintained by computers all over the world. The users of the hosts compose queries, and the search engine identifies pages that match the queries, e.g., pages that include key words of the queries. These pages are known as a result set.\nIn many cases, particularly when a query is short or not well defined, the result set can be quite large, for example, thousands of pages. The pages in the result set may or may not satisfy the user's actual information needs. Therefore, techniques have been developed to identify a smaller set of related pages.\nIn one prior art technique used by the Excite search engine, please see \"http://www.excite.com,\" users first form a query that attempts to specify a topic of interest. After the result set has been returned, the user can use a \"Find Similar\" option to locate related pages. However, there the finding of the related pages is not fully automatic because the user first is required to form a query, before related pages can be identified. In addition, this technique only works on the Excite search engine and for the specific subset of Web pages that are indexed by the Excite search engine.\nIn another prior art technique, an algorithm for connectivity analysis of a neighborhood graph (n-graph) is described by Kleinberg in \"Authoritative Sources in a Hyperlinked Environment,\" Proc. 9th ACM-SIAM Symposium on Discrete Algorithms, 1998, and also in IBM Research Report RJ 10076, May 1997, see, \"http://www.cs .cornell.edu/Info/People/kleinber/auth.ps.\" The algorithm analyzes the link structure, or connectivity of Web pages \"in the vicinity\" of the result set to suggest useful pages in the context of the search that was performed.\nThe vicinity of a Web page is defined by the hyperlinks that connect the page to others. A Web page can point to other pages, and the page can be pointed to by other pages. Close pages are directly linked, farther pages are indirectly linked. This connectivity can be expressed as a graph where nodes represent the pages, and the directed edges represent the links. The vicinity of all the pages in the result set is called the neighborhood graph.\nSpecifically, the Kleinberg algorithm attempts to identify \"hub\" and \"authority\" pages in the neighborhood graph for a user query. Hubs and authorities exhibit a mutually reinforcing relationship.\nIn U.S. patent application Ser. No. 09/007,635 \"Method for Ranking Pages Using Connectivity and Content Analysis\" filed by Bharat et al. on Jan. 15, 1998, a method is described that examines both the connectivity and the content of pages to identify useful pages. However, the method is relatively slow because all pages in the neighborhood graph are fetched in order to determine their relevance to the query topic. This is necessary to reduce the effect of non-relevant pages in the subsequent connectivity analysis phase.\nIn U.S. patent application Ser. No. 09/058,577 \"Method for Ranking Documents in a Hyperlinked Environment using Connectivity and Selective Content Analysis\" filed by Bharat et al. on Apr. 9, 1998, a method is described which performs content analysis only a small subset of the pages in the neighborhood graph to determine relevance weights, and pages with low relevance weights are pruned from the graph. Then, the pruned graphed is ranked according to a connectivity analysis. This method still requires the result set of a query to form a query topic.\nIn any of the above cases, it would be advantageous if duplicate or near duplicate pages could quickly be identified since these pages essentially represent the same content. It would even be better, if near duplicates could be identified without having the analyze the detailed content of the pages."}
-{"text": "1. Field of the Invention\nThe invention relates to a method of manufacturing a plurality of electronic multilayer components, each of which comprises alternately stacked electrically conductive and electrically insulating layers, the electrically conductive layers being electrically connected in a periodically alternate arrangement to different edges of the multilayer component. Such components may receive application as multilayer capacitors or multilayer actuators, example.\n2. Discussion of the Related Art\nA method as described in the opening paragraph is known from U.S. Pat. No. 3,326,718, in which layers of electrically conductive and electrically insulating material are alternately deposited onto a flat substrate through an apertured mask, the planes of the substrate and mask being mutually parallel. In the case of the insulating material, the depository flux is directed at right angles to the plane of the mask, so that it passes through the aperture in a perpendicular direction. However, in the case of the conductive material, the depository flux is directed through the aperture at a non-perpendicular angle .alpha. with respect to the substrate surface. Moreover, although consecutive conductive layers are deposited using the same value of .alpha., the depository fluxes for such consecutive layers are not mutually parallel, but instead arise from sources located at diametrically opposite sides of the aperture. As a result, consecutive conductive layers demonstrate only a partial mutual overlap, as illustrated in FIG. 2 of the cited U.S. patent. At the same time, as shown in FIG. 4 of that patent, conductive layers (56, 56') having an odd ordinal number make mutual electrical contact at one side (62) of the component, and conductive layers (68, 68') having an even ordinal number make mutual electrical contact at an other side (76) of the component.\nThe known method has a number of disadvantages. In particular, the number of conductive layers which can be deposited in this manner is severely limited. This is because, as the stack of layers on the substrate increases in height, that stack will itself begin to partially eclipse the depository fluxes of conductive material, and will eventually prevent the desired mutual contact of every second conductive layer at the edge of the component. In such a scenario, the finished component will have to be provided along its sides with blanketing layers of conductive material (such as solder), for the purpose of achieving uninterrupted interconnection of the conductive layers terminating at each given side. However, because such blanketing layers are at the sides of the component, they are not directly compatible with surface mounting techniques, which require the component's electrical contacts to be located in a single plane."}
-{"text": "The present invention relates to a universal mechanical control system for moving a controlled device in response to the movement of a universally movable control rod. The invention is particularly applicable for controlling a laser beam surgical scalpel in accordance with the movement of a manipulatable joy stick and is therefore described below with respect to this application.\nLaser beam surgical scalpels are increasingly being used for performing delicate surgical operations. A known type of such scalpels includes a microscope to permit viewing by the surgeon of the working area, and an optical control unit including a control rod, commonly called a joy stick, manipulatable by the surgeon for directing the laser beam with respect to the working area. Critical factors in the successful use of such scalpels for performing delicate surgery are the ease and precision by which the laser beam may be manipulated by the surgeon by using the joy stick. Various types of mechanical and electrical control systems have been proposed or are now in use for this purpose, but the known systems are generally of complicated, bulky, and expensive construction and/or are not entirely satisfactory with respect to the ease and precision by which they may be conveniently manipulated by the surgeon during the performance of a delicate surgical procedure.\nAn object of the present invention is to provide a simple, inexpensive and easily-operated mechanical control system for moving a controlled device in response to the movement of a universally movable control member. Another object of the invention is to provide a laser beam surgical scalpel with the above mechanical control system for manipulating the laser beam with respect to the working area in response to the manipulations of a joy stick."}
-{"text": "Relays may be used in LTE-A to enhance coverage and capacity and offer more flexible deployment options. A relay node (RN) may create new cells, distinguishable and separate from the cells of a donor-evolved NodeB (eNB), where the donor-eNB may support a Un interface to support eNB to RN communication. To any legacy wireless transmit/receive unit (WTRU), an RN may appear as an eNB. That is, the presence of an RN in its communication path to the donor-eNB may be transparent to the WTRU. An RN may be an eNB that has a wireless in-band backhaul link back to the donor-eNB by using an LTE or LTE-A air interface within the International Mobile Telecommunications (IMT) spectrum allocation.\nA two-way communications system with two transceivers and one relay node may take four time slots to complete a message exchange, assuming time division duplex (TDD) half duplex mode and no direct link between two transceivers. It would be desirable to have a method and apparatus for implementing a relay scheme that may use a lower number of time slots and may also have a low level of complexity."}
-{"text": "Caps and hats vary greatly in terms of design and functionality. While baseball hats are very trendy in contemporary times, they fail to provide the warmth necessary for extreme drops in temperature. A typical baseball cap is constructed of a webbed plastic mesh. Such a hat prevents sun from obstructing the wearer's view, but the hat does little in the way of maintaining body heat. In fact, baseball hats are generally worn as ornamental pieces outside the realm of sporting events. Although some baseball hats are constructed of heavy cotton and other materials, the basic skeleton of a baseball hat forces the hat to sit atop the wearer's head, not on and around the wearer's head.\nSnow hats, on the other hand, are traditionally constructed primarily of wool and acrylic. They are designed to keep the wearer's head warm. In this regard, a typical snow hat exhibits uniform construction and has little structural integrity when removed from the wearer's head. Snow hats are not intended to be rigged and present obstructions to the sun, wind, and rain; but rather, snow hats are designed to be close fitting and insulating with respect to the wearer's head, maintaining anterior temperature as moisture and heat are conserved.\nThe prior art is replete with various types of hats, none of which approach the design and functionality of the present invention.\nU.S. Pat. No. 80,352, issued to Ibach and Weidenman on Jul. 28, 1868, illustrates a method of attaching a paper visor to a paper cap wherein a visor and a loop of material is fitted to the bottom rim of a skullcap. When the cap becomes dirty, the visor can be attached to another cap, whereas in the present invention, the visor must firmly remain in combination with a head covering to remain fully effective. Moreover, unlike the present invention, Ibach and Weidenman's device does not flexibly surround the user's head to promote heat retention. U.S. Pat. No. 442,921, issued to Stohr on Dec. 16, 1890, shows a knitted cap with a visor in which a hat with depending ear flap and ties is combined with a visor. Stohr, unlike the present invention, does not provide any means of capturing moisture from the brow or forehead of the user during athletic activity. Furthermore, Stohr's invention, unlike the present invention, employs incisions in the cap to attach the visor.\nU.S. Pat. No. 2,149,655, issued to Yamaguchi on Aug. 3, 1938, illustrates a head covering made entirely of crocheting yarn. The device discloses a means for reinforcing a crocheted visor to a crocheted hat. Unlike the present invention, Yamaguchi's device is concerned with forming a bead at the junction of the crown and the visor to add stability to stiching.\nU.S. Pat. No. 2,158,861, issued to Meyer on Sep. 2, 1937, shows a visor for collapsible caps which is capable of assuming a flat position when a collapsible cap is in a folded, flat position. Unlike the present invention, Meyer's device is not capable of providing a visor in combination with a knit type hat.\nU.S. Pat. No. 2,417,986, issued to Marder et al. on Mar. 25, 1947, is directed to a cap and visor with draw strings extending to the nape of the user's neck. Unlike the present invention, Marder's invention has multiple lines of stitching and is secured behind the user's head.\nU.S. Pat. No. 2,420,569, issued to Sewell on May 13, 1947, shows a baseball cap. Unlike the present invention, Sewell's device has a visor which is attached to the crown of the cap.\nU.S. Pat. No. 2,651,044, issued to Stankiewicz et al. on Sep. 8, 1953, depicts a scarf hat having interfitting concentric outer and inner head band members for receiving between them an edge of a scarf. Unlike the present invention, the visor is attached by fitting between the outer and inner head band members.\nU.S. Pat. No. 4,601,070, issued to Sargentini on Jul. 22, 1986, shows a novelty ski hat. Unlike the present invention, there is no visor attached to the hat.\nU.S. Pat. No. 5,471,684, issued to Casale on Dec. 5, 1995, shows a sports cap with a sliding brim. Unlike the present invention, the brim is detachably secured to the bottom edge of the cap. Furthermore, unlike the present invention, the majority of the cap can be detached from the bottom edge of the cap.\nU.S. Pat. No. 5,481,759, issued to Rinaldi on Jan. 9, 1996, illustrates an expandable baseball hat and cover. Unlike the present invention, Rinaldi's device employs tabs engaged to gather flexible material for a snug fit. The present invention, by nature of its design, fits snugly to the user's head without employing a complicated tab system.\nFrench Publication No. 621,305, invented by Menant and published on May 9, 1927, shows a woven cap. Unlike the present invention, Menant's device employs multiple stiches about the middle periphery of the cap. Further, Menant's device is not designed, by virtue of its shape and inability to snugly fit the user's head, for athletic use.\nAccordingly, the need arises for a multi functional hat with a brim or visor, which is a blends the functionality of both the typical snow hat and the typical baseball cap. The multi functional hat must be rigid so as to protect the wearer's face from the outdoor elements. Moreover, the multi functional hat should be capable of maintaining the warmth of the wearer's head despite baseball cap type adaptations. The multi functional hat should be of a design that does not compromise the advantages of existing snow hats, while at the same time, exploiting the current weaknesses of headgear design."}
-{"text": "The present invention relates to an improved trimmer head for use in flexible line rotary trimming devices used to trim grass, weeds and other vegetation. More particularly, the invention is directed to an improvement in the line loading of both \u201cbump-feed\u201d type trimmer heads such as those disclosed in U.S. Pat Nos. 4,458,419 and 4,959,904 and the more simple manually operated heads such as that disclosed in U.S. Pat. No. 4,145,809, the contents of said patents being incorporated herein by reference as though fully set forth below.\nTrimmer heads used in flexible line rotary trimmers generally carry one or two lengths of flexible nylon cutting line wrapped about an interior spool with the ends of the line or lines projecting outwardly through opposed apertures in the side wall of the trimmer head. The head is threadably mounted on the end of an elongated shaft and rotated at a high velocity by a gas or electric motor so that the ends of the cutting line project radially from the head and sever weeds or other vegetation. When cutting line projecting from the head breaks off or becomes overly worn, it must be severed and fresh line extended from the spool through the line outlet eyelets in the side of the housing. Bump-feed type heads include a line feed-out mechanism which responds to a bump on the ground intentionally applied by the operator to feed out a measured length of fresh cutting line which is typically cut to the desired length by a knife blade projecting from a shield attached to the trimmer above the cutting head and spaced a predetermined distance from the perimeter of the trimmer head housing. Manual heads do not include any such line feed-out mechanism. A fastening nut that holds the housing portion of the trimmer head to the spool must be loosened so that the spool can be separated from the housing and manually rotated relative to the housing to pay out additional cutting line. The spool and housing are then re-secured by the fastening member.\nIn both bump-feed and manual heads, the length or lengths of cutting line are typically wound onto the spool by hand. As most cutting heads employ two lengths of line projecting from opposed sides of the cutting head, care must be taken during the winding of the spool to avoid crossing or otherwise tangling of the two lines within the spool which interferes with the paying out of fresh line. This is particularly important in bump-feed heads where centrifugal force is utilized to pull the new lengths of line from the spool during use as the head is being bumped against the ground as any line tangle will interfere with the proper feeding of the line. Difficulty in properly loading the line on the spool is the most common complaint of home users of flexible line trimmers. It also is a time consuming task for the professional user.\nThe early bump-feed mechanisms typically consisted of a dog or friction clutch located between the spool of line and the surrounding housing. By bumping an extension of the spool on the ground, or other fixed object, the friction clutch was temporarily disengaged for a length of time dependant on the duration of the bump. The dog clutch released by the bump then abruptly engaged at the next opportunity to feed out line in segment lengths which were related to the engagement points of the dog clutch. Such dog clutches had outwardly extending ribs which engaged inwardly extending abutment tangs and therefore depended upon a skillful bump when it was desired to deliver only one segment length. However, friction within such devices and overzealous bumping often resulted in two or more line segments being fed out, particularly especially when the device has been in use and the corners on the ribs and tangs became worn such that positive engagement was no longer assured. The unavoidable abrupt operation of the dog clutch caused such wear to occur.\nA bump-feed-out mechanism was subsequently developed that automatically fed out a predetermined length of line with each bump, regardless of the duration of the bump, and which did not lose this ability with extended use. That device is disclosed in and is the subject of the incorporated reference, U.S. Pat. No. 4,458,419. As described therein in detail, the improved trimmer head contained a spool holding one or more coils of cutting line and a simplified mechanism that selectively allowed relative movement of the spool with respect to the housing in response to bumping of the head on the ground to pay out measured lengths of line. The simplified pay-out mechanism included a novel spring-loaded cam and cam follower arrangement in which the cam follower included two pair of diametrically opposed and generally inwardly facing abutment surfaces arrayed about the axis of rotation of the trimmer housing. The abutment surfaces were thus spaced 90\u00b0 apart and were carried by a depending cylindrical wall that circumscribed an interior chamber. The cam member was disposed within the chamber in threaded engagement with the extended lower end of the drive bolt of the trimmer housing and defined two vertically adjacent cams, each cam being of a square configuration and defining four perpendicularly disposed cam surfaces adapted to engage the abutment surfaces on the cam follower. The upper cam was rotationally offset 45\u00b0 from the lower cam.\nIn operation, the housing was rotationally driven by the drive bolt through a connection between the upper end of the bolt and the trimmer drive means. The housing and cam member was thus driven by the drive bolt, which in turn drove the cam follower and the spool mounted thereon due to the engagement between the cam surfaces on the cam member and the abutment surfaces on the cam follower. The line carrying spool was disposed about the cylindrical wall of the cam follower and attached thereto via a pair of opposed outwardly projecting studs on the cam follower member that extend into slots formed in the inner portion of the spool. The spool was provided with a bumper at its lower end such that when the bumper was pressed against or bumped on the ground, the housing moved downwardly with respect to the spool against the force of a spring, disengaged the lower cam from the abutment surfaces on the cam follower and allowed the cam member to rotate 45\u00b0 relative to the cam follower, whereupon the cam surfaces of the upper cam would abut the abutment surfaces on the cam follower. That imparted a similar degree of relative rotation between the spool and the housing. Once the force of the bump was dissipated, the spring loading forced the spool and housing back to their relative positions, which released the cam surfaces on the upper cam from the cam follower abutment surfaces and allowed another 45\u00b0 of relative rotation of the cam member and cam follower and thus of the spool and housing, for a total of 90\u00b0 of rotation per bump, which provided the predetermined relative rotation between the housing and spool needed to pay out a desired length of line through the apertures in the trimmer housing. Since the cams interacted with simple, inwardly facing cam follower surfaces formed only on a single level, the release mechanism was deemed relatively economical to manufacture and, due to the large abutment surface areas presented between the cams and cam follower, the device was durable, trouble free and reliable.\nBecause of early difficulties in molding some of the components of the cutting head disclosed in U.S. Pat. No. 4,458,419, the head became more expensive to manufacture than anticipated. New material developments subsequently reduced the cost of manufacture. In the meantime, however, a similar bump-feed drive mechanism was developed in which large square cams were formed on the upper and lower outer radial surfaces of the spool and the corresponding cam followers were formed by upper and lower portions of the housing which surrounded the spool. Such a head is disclosed in the incorporated reference, U.S. Pat. No. 4,959,904, and is still in production.\nOver the years, with increasing competition from offshore manufacturers, it became clear that even with the development of new materials the earlier bump-feed mechanism covered by U.S. Pat. No. 4,459,419 was not as economical to manufacture as earlier believed. It contained several parts, some of which had to be hand assembled. In addition, vibration, the threaded engagement between the cam member and the drive bolt, and the heat generated by the trimmer required the use of a chemical bonding agent having a high melting point to prevent the cam member and cam follower from breaking loose from the drive bolt. Such agents, however, had extremely high break way torques, rendering the threaded connection effectively permanent. As a result, certain components of the head could not be replaced when worn. Thus, that head was significantly modified so as to retain all of the advantages of its predecessor yet utilize fewer component parts and obviate the need for any hand assembly and use of chemical bonding. That modification is the subject of a pending U.S. patent application entitled \u201cTrimmer Head for Use in Flexible Line Rotary Trimmers\u201d, filed Oct. 2, 2003 and identified by Ser. No. 10/677,700, of which this application is a continuation-in-part. The bump-feed and manual heads of the present invention retain all of the advantages of the above-described heads and adds thereto the ability to far more quickly and easily uniformly wind lengths of cutting line onto the spool without materially increasing the cost of production."}
-{"text": "S-Adenosylmethionine (SAMe) is found in almost every tissue and fluid in the body. SAM plays a crucial role in the process called transmethylation. Methylation is involved in nearly every aspect of life. SAM is the primary \u201cmethyl\u201d donor for a variety of methyl-transfer reactions in DNA, RNA, proteins, lipids, and small molecules in the body. Proper DNA methylation is essential for normal embryonic development. Methyl-transferase gene homozygously deleted (knocked out) has been proven lethal (Pegg, A. E., Feith, D. J., Fong, L. Y., Coleman, C. S., O'Brian, T. G., and Shantz, L. M., 2003, Biochem. Soc. Trans. 31, 356-360). DNA improperly methylated has been found in many tumors. Alterations in DNA methylation patterns induce the expression of oncogens or silence the expression of tumor suppressor genes, and methyl deficient diets have been shown to promote liver cancer in rodents.\nThe transsulfuration begins with S-adenosylhomocysteine (SAH), the residual structure of SAM upon donating the methyl group (transmethylation). Hydrolysis of SAH yields homocysteine, which in turns converts to cystathionine, then cysteine, and eventually, to glutathione, the hepatocellular antioxidant and life-saving detoxification agent.\nThe aminopropylation is another process initiated with SAM through decarboxylation. The decarboxylated SAM then couples with putrescine to generate spermidine and spermine which are critical to cell growth, differentiation and the stability of DNA and RNA. Furthermore, Methylthioadenosine (MTA), the by-product of polyamine synthesis, is a powerful analgesic and anti-inflammatory agent. This may be, at least partially, responsible for the clinical benefits observed in the treatment of osteoarthritis, rheumatoid arthritis and fibromyalgia with SAMe.\nSAMe plays a role in the immune system, maintains cell membranes, and helps produce and break down brain chemicals, such as serotonin, melatonin, and dopamine. Deficiency of either vitamin B12 or foliate can reduce the level of SAMe. SAMe is also an antioxidant, a substance that protects the body from damaging reactive oxygen molecules in the body. These reactive oxygen molecules can come from inside the body or from environmental pollution and are thought to play a role in the aging process and the development of degenerative disease. In general, SAMe is thought to raise the level of functioning of other amino acids in the body.\nBy way of further background, S-adenosyl-1-methionine is a substrate of an enzyme lyase that converts S-adenosyl-1-methionine to the molecule methylthioadenosine and homoserine; it is an aminobutyric chain donor to tRNA; it is an aminoacidic chain donor in the biosynthesis of biotin; SAM-e, after decarboxylation, is the donor of aminopropyl groups for the biosynthesis of neuroregulatory polyamines spermidine and spermine. (Zappia et al (1979), Biomedical and Pharmacologcial roles of Adenosylmethionine and the Central Nervous System, page 1, Pergamon Press. N.Y.)\nSAM-e has been used clinically in the treatment of liver disease (Friedel H, Goa, K. L., and Benfield P., (1989), S-Adenosyl-1-methionine: a review of its pharmacological properties and therapeutic potential in liver dysfunction and affective disorders in relation to its physiological role in cell metabolism. Drugs. 38, 389-416), arthritis (Di Padova C, (1987), S-adenosyl-1-methionine in the treatment of osteoarthritis: review of the clinical studies. Am J. Med. 83, (Suppl. 5), 6-65), and depression (Kagan, B, Sultzer D. L., Rosenlicht N and Gerner R. (1990), Oral S-adenosylmethionine in depression: a randomized, double blind, placebo-controlled trial. Am. J. Psychiatry 147, 591-595.) Alzheimer's patients have reduced cerebral spinal fluid levels of S-adenosyl-1-methionine (Bottiglieri et al, (1990), Cerebrospinal fluid S-adenosyl-1-methionine in depression and dementia: effects of treatment with parenteral and oral 5-adenosyl-1-methionine. J. Neurol. Neurosurg. Psychiatry 53, 1096-1098.) In a preliminary study, SAM-e was able to produce cognitive improvement in patients with Alzheimer's disease. (Bottiglieri et al (1994), The clinical potential of admetionine (S-adenosyl-1-methioinine) in neurological disorders. Drugs 48, 137-152.) SAM-e brain levels in patients with Alzheimer's disease are also severely decreased. (Morrison et al, (1996), Brain S-adenosylmethionine levels are severely decreased in Alzheimer's disease, Journal of Neurochemistry, 67, 1328-1331.) Patients with Parkinson's disease have also been shown to have significantly decreased blood levels of SAM-e. (Cheng et al, (1997), Levels of L-methionine S-adenosyltransferase activity in erythrocytes and concentrations of S-adenosylmethionine and S-adenosylhomocysteine in whole blood of patients with Parkinson's disease. Experimental Neurology 145, 580-585.)\nSAM-e levels in patients treated with the antineoplastic drug methotrexate are reduced. Neurotoxicity associated with this drug may be attenuated by co-administration of SAM-e. (Bottiglieri et al (1994), The Clinical Potential of Ademetionine (S-adenosylmethionine) in neurological disorders, Drugs, 48 (2), 137-152.)\nCerebral spinal fluid levels of SAM-e have been investigated in HIV AIDS dementia Complex/HIV encephalopathy and found to be significantly lower than in non-HIV infected patients. (Keating et al (1991), Evidence of brain methyltransferase inhibition and early brain involvement in HIV positive patients Lancet: 337:935-9.)\nDe La Cruz et al have shown that SAM-e, chronically administered, can modify the oxidative status in the brain by enhancing anti-oxidative defenses. (De La Cruz et al, (2000), Effects of chronic administration of S-adenosyl-1-methionine on brain oxidative stress in rats. Naunyn-Schmiedeberg's Archives Pharmacol 361: 47-52.) This is similar to results obtained with SAM-e in liver and kidney tissue. Thus SAM-e would be useful as an antioxidant.\nOral SAM-e administration to patients with and without liver disease has resulted in increases in liver glutathione levels. (Vendemiale G et al, (1989), Effect of oral S-adenosyl-1-methionine on hepatic glutathione in patients with liver disease. Scand J Gastroenterol; 24: 407-15. Oral administration of SAM-e to patients suffering from intrahepatic cholestasis had improvements in both the pruritus as well as the biochemical markers of cholestasis. (Giudici et al, The use of admethionine (SAM-e) in the treatment of cholestatic liver disorders. Meta-analysis of clinical trials. In: Mato et al editors. Methionine Metabolism: Molecular Mechanism and Clinical Implications. Madrid: CSIC Press; 1992 pp 67-79.) Oral SAM-e administration to patients suffering from primary fibromyalgia resulted in significant improvement after a short term trial. (Tavoni et al, Evaluation of S-adenosylmethioine in Primary Fibromaylgia. The American Journal of Medicine, Vol 83 (suppl 5A), pp 107-110, 1987.) SAM-e has been used for the treatment of osteoarthritis as well. (Koenig B. A long-term (two years) clinical trial with S-adenosylmethionine for the treatment of osteoarthritis. The American Journal of Medicine, Vol 83 (suppl 5A), Nov. 20, 1987 pp 89-94)\nSAM-e is clinically useful in many apparently unrelated areas because of its important function in basic metabolic processes. One of its most striking clinical uses is in the treatment of alcoholic liver cirrhosis that, until now, remained medically untreatable. Mato et al demonstrated the ability of oral SAM-e in alcoholic liver cirrhosis to decrease the overall mortality and/or progression to liver transplant by 29% vs 12% as compared with a placebo treated group. (Mato et al (1999), S-adenosylmethionine in alcohol liver cirrhosis: a randomized, placebo-controlled, double blind, multi-center clinical trial, Journal of Hepatology, 30, 1081-1089.)\nSam-e also attenuates the damage caused by tumor necrosis factor alpha and can also decrease the amount of tumor necrosis factor alpha secreted by cells. Consequently, conditions in which this particular inflammatory factor is elevated would benefit from the administration of SAM-e. (Watson W H, Zhao Y, Chawla R K, (1999) Biochem J August 15; 342 (Pt 1):21-5. S-adenosylmethionine attenuates the lipopolysaccharide-induced expression of the gene for tumour necrosis factor alpha.) SAM-e has also been studied for its ability to reduce the toxicity associated with administration of cyclosporine A, a powerful immunosuppressor. (Galan A, et al, Cyclosporine A toxicity and effect of the s-adenosylmethionine, Ars Pharmaceutica, 40:3; 151-163, 1999.)\nSAM-e, incubated in vitro with human erythrocytes, penetrates the cell membrane and increases ATP within the cell thus restoring the cell shape. (Friedel et al, S-adenosyl-1-methionine: A review of its pharmacological properties and therapeutic potential in liver dysfunction and affective disorders in relation to its physiological role in cell metabolism, Drugs 38 (3):389-416, 1989)\nSAM-e has been studied in patients suffering from migraines and found to be of benefit. (Friedel et al, S-adenosyl-1-methionine: A review of its pharmacological properties and therapeutic potential in liver dysfunction and affective disorders in relation to its physiological role in cell metabolism, Drugs 38 (3): 389-416, 1989)\nSAM-e has been administered to patients with peripheral occlusive arterial disease and was shown to reduce blood viscosity, chiefly via its effect on erythrocyte deformability.\nSAM-e is commercially available using fermentation technologies that result in SAM-e formulations varying between 60 and 80% purity. (That is, the final product contains 60-80% of the active or (S, S)-SAM-e and 20-40% of the inactive or (R, S)-SAM-e.) (Gross, A., Geresh, S., and Whitesides, Gm (1983) Appl. Biochem. Biotech. 8, 415.) Enzymatic synthetic methodologies have been reported to yield the inactive isomer in concentrations exceeding 60%. (Matos, J R, Rauschel F M, Wong, C H. S-Adenosylmethionine: Studies on Chemical and Enzymatic Synthesis. Biotechnology and Applied Biochemistry 9, 39-52 (1987). Enantiomeric separation technologies have been reported to resolve the pure active enantiomer of SAM-e. (Matos, J R, Rauschel F M, Wong, C H. S-Adenosylmethionine: Studies on Chemical and Enzymatic Synthesis. Biotechnology and Applied Biochemistry 9, 39-52 (1987; Hoffman, Chromatographic Analysis of the Chiral and Covalent Instability of S-adenosyl-1-methionine, Biochemistry 1986, 25 4444-4449: Segal D and Eichler D, The Specificity of Interaction between S-adenosyl-1-methionine and a nucleolar 2-0-methyltransferase, Archives of Biochemistry and Biophysics, Vol. 275, No. 2, December, pp. 334-343, 1989) Newer separation technologies exist to resolve enantiomers on a large commercial production scale at a very economic cost. In addition, it would be conceivable to synthesize the biologically active enantiomer using special sterioselective methodologies but this has not been accomplished to date.\nDe la Haba first showed that the sulfur is chiral and that only one of the two possible configurations was synthesized and used biologically. (De la Haba et al J. Am. Chem. Soc. 81, 3975-3980, 1959) Methylation of RNA and DNA is essential for normal cellular growth. This methylation is carried out using SAM-e as the sole or major methyl donor with the reaction being carried out by a methyltransferase enzyme. Segal and Eichler showed that the enzyme bound (S, S)-SAM-e 10 fold more tightly than the biologically inactive (R, S)-SAM-e thus demonstrating a novel binding stereospecificity at the sulfur chiral center. Other methyltransferases have been reported to bind (R, S)-SAM-e to the same extent as (S, S)-SAM-e and thus (R, S)-SAM-e could act as a competitive inhibitor of that enzyme. (Segal D and Eichler D, The Specificity of Interaction between S-adenosyl-1-methionine and a nucleolar 2-0-methyltransferase, Archives of Biochemistry and Biophysics, Vol. 275, No. 2, December pp. 334-343, 1989; Borchardt R T and Wu Y S, Potential inhibitors of S-adenosylmethionine-dependent methyltransferases. Role of the Asymmetric Sulfonium Pole in the Enzymatic binding of S-adenosyl-1-methionine, Journal of Medicinal Chemistry, 1976, Vol 19, No. 9, 1099-1103.)\nSAM-e (whether in its optically pure enantiomeric form or in an enantiomeric or racemic mixture) presents certain difficult problems in terms of its stability at ambient temperature that result in degradation of the molecule to undesirable degradation products. SAM-e (and thus its enantiomers) must be further stabilized since it exhibits intramolecular instability that causes the destabilization and breakdown of the molecule at both high as well as ambient temperatures. SAM-e has therefore been the subject of many patents directed both towards obtaining new stable salts, and towards the provision of preparation processes that can be implemented on an industrial scale. The present patent thus envisions the use of any of the salts of SAM-e already disclosed in the prior art to stabilize the enantiomeric forms of SAM-e.\nThe clinical diagnostic field has seen a broad expansion in recent years, both as to the variety of materials of interest that may be readily and accurately determined, as well as the methods for the determination. Over the last several decades, testing for numerous substances such as drugs of abuse, or other biological molecules of interest has become commonplace. In recent years, immunoassay based on the interaction of an antibody with an antigen has been extensively investigated for this purpose. Based on the unique specificity and high affinity of antibodies, an immunoassay can accurately and precisely quantitate substances at the very low concentrations found in biological fluids.\nAccordingly, there is a need for improved methods of detection and diagnosis of cancer and other diseases, as well as methods for monitoring the progress of the diseases and monitoring the progress of various treatments for cancer and other diseases by quantitating the methylating index as a biomarker."}
-{"text": "The present invention relates to the structure of a casing of a miniature portable apparatus and a clip attachable thereto and, more particularly, to a structure which allows a clip to be firmly affixed to the casing of a radio pager while promoting easy attachment and detachment of the clip.\nA wide variety of miniature portable apparatuses are extensively use today and include radio pagers and audio apparatus. A radio pager, for example, may be provided with a clip to be put in the user's chest pocket or on the user's waist belt. The clip is attached to the casing of such a radio pager. Specifically, the casing is usually made of plastic and provided with an elongate groove on the top thereof and an elongate lug on the bottom, while an elastic bate plate is made of metal and provided with an engaging portion engageable with the groove in a upper end portion thereof and a slot engageable with the lug in a generally J-shaped lower end portion. The clip has a clip body which is rotatably supported by the base plate. To attach the clip to the casing, the engaging portion of the base plate is put in the groove of the casing, and then the base plate is rotated about the groove. As a result, the slot of the base plate mates with the lug of the casing due to the resiliency of the base plate. To remove the clip from the casing, the lower end portion of the base plate is urged away from the casing by, for example, a tool to release the slot from the lug.\nHowever, the problem with the conventional clip structure is that the edge of the lower end of base plate moves in contact with the casing in the event when the clip is attached to or detached from the casing, scratching or even breaking the casing. Another problem is that the resiliency of the base plate used to attach and detach the clip from the casing is apt to differ from one base plate to another for production reasons. Therefore, some base plate come off the casing rather easy, while other are difficult to remove from the casing. In addition, since the edge of the base plate is forcibly urged away from the casing by a tool, it is likely that the casing is scratched or the base plate is deformed."}
-{"text": "1. Field of the Invention\nThe present invention relates to an exposure technique for exposing a substrate via a mask and a projection optical system, which is suitable for exposing a long sheet-shaped photosensitive object wound in a roll shape. In addition, the present invention relates to a fabricating technique for fabricating (manufacturing) a device using such an exposure technique.\n2. Description of Related Art\nFor example, in an exposure apparatus that is used for fabricating an element such as a semiconductor element or a liquid crystal display element, a typical exposure target (exposure object) has conventionally been a flat plate-shaped object which has high rigidity, such as a glass substrate or a semiconductor wafer which have been coated with photoresist. In recent years, in order to efficiently fabricate a device having a large area, as an exposure target, a long sheet-shaped member which is flexible and can be stored by being wound up into a roll shape has been used. When exposing such a long sheet-shaped member, as a conventional example, a part of the sheet-shaped member is exposed via a mask and a projection optical system is subjected to one-shot exposure (still exposure), the sheet-shaped member is then suctioned and moved, for example, by an exposure table, and then the suction is released and only the exposure table is returned to the initial position. The above described operation is repeatedly carried out to intermittently move the sheet-shaped member from a supply roller to a winding-up roller (for example, see Japanese Patent Application Publication No. 2007-114385).\nAccording to the conventional exposure method of the sheet-shaped member, one-shot exposure via a mask and movement of the sheet-shaped member are alternately repeated. Therefore, exposure efficiency is low, and, for example, it takes a long time to expose a sheet-shaped member which is wound up into one roll.\nIn contrast, a method is also considered in which the sheet-shaped member is scan-exposed while a mask and the sheet-shaped member are moved synchronously. However, in the case of mere scan exposure, for example, if the sheet-shaped member is exposed in an outward path of the mask, then it is difficult to expose the sheet-shaped member in a return path of the mask. For this reason, substantially half of the reciprocation time of the mask becomes a period at which exposure is impossible, and exposure efficiency cannot be high."}
-{"text": "1. Field of the Invention\nThe present invention relates to a delay circuit, a testing apparatus, and a capacitor. More particularly, the present invention relates to a delay circuit which generates a desired delay time by changing the junction capacitance of a field effect transistor.\n2. Description of the Related Art\nFIG. 1 shows a conventional delay circuit 300. The conventional delay circuit 300 has a first buffer 302 which shapes the wave form of an input signal and then outputs the resultant shaped signal, a path 306 through which the output signal transmits, a first capacitor 312 which adds capacitance C to the path 306, a second capacitor 314 which adds capacitance C\u2032 to the path 306, a first switching device 308 which electrically connects or disconnects the path 306 with the first capacitor 312, a second switching device 310 which electrically connects or disconnects the path 306 with the second capacitor 314, and a second buffer 304 which shapes the wave form of the signal that has transmitted through the path 306 and outputs the resultant shaped signal. A control unit not shown in the drawing controls the switching devices 308 and 310 so as to change the capacitance added to the path 306. In this way, the control unit not shown in the drawing delays the signal that transmits the path 306 by a desired length of time.\nThe conventional delay circuit 300 achieves a fine delay resolution by selectively adding either the capacitance C or the capacitance C\u2032 which differs slightly from the capacitance C. However, in the conventional delay circuit 300, the channel capacitance of the first switching device 308 differs from that of the second switching device 310, and the wire capacitance of the wire which connects the first capacitor 312 with the path 306 differs from the wire capacitance of the wire which connects the second capacitor 314 with the path 306. These capacitance differences influence the capacitance added to the path 306. As a result, the desired fine delay resolution which is designed to be achieved by utilizing the fine difference between the capacitance C and the capacitance C\u2032 has been very difficult."}
-{"text": "Vehicle windows (e.g., windshields, backlites, sunroofs, and sidelites) are known in the art. For purposes of example, vehicle windshields typically include a pair of bent glass substrates laminated together via a polymer interlayer such as polyvinyl butyral (PVB). It is known that one of the two glass substrates may have a coating (e.g., low-E coating) thereon for solar control purposes such as reflecting R and/or UV radiation, so that the vehicle interior can be more comfortable in certain weather conditions. Conventional vehicle windshields are made as follows. First and second flat glass substrates are provided, one of them optionally having a low-E coating sputtered thereon. The pair of glass substrates are washed and booked together (i.e., stacked on one another), and then while booked are heat bent together into the desired windshield shape at a high temperature(s) (e.g., 8 minutes at about 600-625 degrees C.). The two bent glass substrates are then laminated together via the polymer interlayer to form the vehicle windshield.\nInsulating glass (IG) window units are also known in the art. Conventional IG window units include at least first and second glass substrates (one of which may have a solar control coating on an interior surface thereof) that are coupled to one another via at least one seal(s) or spacer(s). The resulting space or gap between the glass substrates may or may not be filled with gas and/or evacuated to a low pressure in different instances. However, many IG units are required to be tempered. Thermal tempering of the glass substrates for such IG units typically requires heating the glass substrates to temperature(s) of at least about 600 degrees C. for a sufficient period of time to enable thermal tempering.\nOther types of coated articles also require heat treatment (HT) (e.g., tempering, heat bending, and/or heat strengthening) in certain applications. For example and without limitation, glass shower doors, glass table tops, and the like require HT in certain instances.\nDiamond-like carbon (DLC) is sometimes known for its scratch resistant properties. For example, different types of DLC are discussed in the following U.S. Pat. Nos. 6,303,226; 6,303,225; 6,261;693; 6,338,901; 6,312,808; 6,280,834; 6,284,377; 6,335,086; 5,858,477; 5,635,245; 5,888,593; 5,135,808; 5,900,342; and 5,470,661, all of which are hereby incorporated herein by reference.\nIt would sometimes be desirable to provide a window unit or other glass article with a protective coating including DLC in order to protect it from scratches and the like. Unfortunately, DLC tends to oxidize and burn off at temperatures of from approximately 380 to 400 degrees C. or higher, as the heat treatment is typically conducted in an atmosphere including oxygen. Thus, it will be appreciated that DLC as a protective overcoat cannot withstand heat treatments (HT) at the extremely high temperatures described above which are often required in the manufacture of vehicle windows, IG window units, glass table tops, and/or the like. Accordingly, DLC cannot be used alone as a coating to be heat treated, because it will oxidize during the heat treatment and substantially disappear as a result of the same (i.e., it will burn off).\nCertain other types of scratch resistant materials also are not capable of withstanding heat treatment sufficient for tempering, heat strengthening and/or bending of an underlying glass substrate.\nAccordingly, those skilled in the art will appreciate that a need in the art exists for a method of making a scratch resistant coated article that is capable of being heat treated (HT) so that after heat treatment the coated article is still scratch resistant. A need for corresponding coated articles, both heat treated and pre-HT, also exists."}
-{"text": "In recent years, an electronic device (compound semiconductor device), in which a GaN layer and an AlGaN layer are disposed sequentially above a substrate and the GaN layer is used as an electron transit layer, has been developed actively. As for an example of such compound semiconductor devices, a GaN based high electron mobility transistor (HEMT) is mentioned. In the GaN based HEMT, a high-concentration two-dimensional electron gas (2DEG) generated at a heterojunction interface between AlGaN and GaN is utilized.\nThe band gap of GaN is 3.4 eV and is larger than the band gap of Si (1.1 eV) and the band gap of GaAs (1.4 eV). That is, GaN has a high breakdown field strength. Furthermore, GaN has also a large saturation electron velocity. Therefore, GaN is a very promising material for a compound semiconductor device capable of high-voltage operation and production of high output, for example, a material for a semiconductor device for a power supply. Consequently, the compound semiconductor device by using the GaN based compound semiconductor device is expected as a high-breakdown voltage power device for a high-efficiency switching element, an electric car, and the like.\nRegarding the GaN based HEMT, the material for a gate electrode is different from the material for the source electrode and the drain electrode. Therefore, the gate electrode is formed by a process different from that of the source electrode and the drain electrode. The gate electrode, the source electrode, and the drain electrode are formed by, for example, a lift-off method. That is, in formation of the electrode, formation of a resist pattern, formation of an electrode material, and removal of the resist pattern are performed. Meanwhile, in production of a GaN based HEMT, recesses or opening portions may be formed in regions to be provided with the gate electrode, the source electrode, and the drain electrode of the compound semiconductor layer. In this case, a compound semiconductor layer is etched by using a resist pattern, so as to form a recess or an opening portion, and thereafter, the resist is removed. Furthermore, after the source electrode and the drain electrode are formed, the source electrode and the drain electrode may be covered with a passivation film, and before formation of the gate electrode, the passivation film may be dry-etched by using a resist pattern. In this case, the resist is removed after the dry etching.\nIn these methods, although the resist is removed, residues of the resist may remain. Moreover, after the dry etching of the passivation film, etching residues may remain. In this regard, these residues may not be removed sufficiently even by cleaning with an organic material. This is because of effects of alteration and the like during post baking after development performed in formation of the resist pattern and dry etching.\nMeanwhile, when the residues are removed by an acid treatment through the use of, for example, a mixture of sulfuric acid and aqueous hydrogen peroxide, an electrode exposed at that time is damaged. For example, in the case where a recess for the gate electrode is formed after the source electrode and the drain electrode are formed, the source electrode and the drain electrode are exposed during formation of the recess. Therefore, if the above-described mixture is used, the source electrode and the drain electrode are damaged. Alternatively, even when the electrodes are not damaged, the exposed surface of the compound semiconductor layer may undergo oxidation, surface roughening, and the like, so that the characteristics may be changed.\nThen, in the case where the above-described residues are present, an increase in leakage current, fluctuations in threshold voltage due to trap of charge, and the like occur, so that the yield is reduced significantly.\nThe followings are reference documents:\n[Document 1] Japanese Laid-open Patent publication No. 2003-167360 and\n[Document 2] Japanese Laid-open Patent publication No. 2007-128038."}
-{"text": "Intracardiac and intravascular procedures, whether performed percutaneously or in an open, surgical, fashion, may liberate particulate debris. Such debris, once free in the vascular system, may cause complications including vascular occlusion, end-organ ischemia, stroke, and heart attack. Ideally, this debris is filtered from the vascular system before it can travel to distal organ beds.\nUsing known filter mechanisms deployed in the arterial system, debris is captured during systole. There is a danger, however, that such debris may escape the filter mechanism during diastole or during filter removal. Apparatus and methods to reduce debris escape during diastole or during filter removal may be desirable to reduce embolic complications."}
-{"text": "The last decade has seen knowledge of the immune system and its regulation expand tremendously. One area of particular interest has been that of research on the proteins and glycoproteins which regulate the immune system. One of the best known families of these molecules are the cytokines. These are molecules which are involved in the \"communication\" of cells with each other. The individual members of the cytokine family have been found to be involved in a wide variety of pathological conditions, such as cancer and allergies. Whereas sometimes the cytokines are involved in the pathology of the condition, they are also known as being therapeutically useful.\nInterleukins are one type of cytokine. The literature on interleukins is vast. An exemplary, but by no means exhaustive listing of the patents in this area includes U.S. Pat. No. 4,778,879 to Mertelsmann et al.; U.S. Pat. No. 4,490,289 to Stern; U.S. Pat. No. 4,518,584 to Market al.; and U.S. Pat. No. 4,851,512 to Miyaji et al., all of which involve interleukin-2 or \"IL-2.\" Additional patents have issued which relate to interleukin-1 (\"IL-1\"), such as U.S. Pat. No. 4,808,611 to Cosman. The disclosure of all of these patents are incorporated by reference herein. More recent patents on different interleukins include U.S. Pat. No. 5,694,234 (IL-13); U.S. Pat. No. 5,650,492 (IL-12); U.S. Pat. Nos. 5,700,664, 5,371,193 and U.S. Pat. No. 5,215,895 (IL-11); U.S. Pat. Nos. 5,728,377, 5,710,251, 5,328,989 (IL-10); U.S. Pat. Nos. 5,580,753, 5,587,302, 5,157,112, 5,208,218 (IL-9); U.S. Pat. Nos. 5,194,375, 4,965,195 (IL-7); U.S. Pat. Nos. 5,723,120, 5,178,856 (IL-6), and U.S. Pat. No. 5,017,691 (IL-4). Even a cursory review of this patent literature shows the diversity of the properties of the members of the interleukin family. One can assume that the larger cytokine family shows even more diversity. See, e.g., Aggarwal et al., ed., Human Cytokines: Handbook For Basic And Clinical Research (Blackwell Scientific Publications, 1992), Paul, ed., Fundamental Immunology (Raven Press, 1993), pg 763-836, \"T-Cell Derived Cytokines And Their Receptors\", and \"Proinflammatory Cytokines and Immunity.\" All cited references are incorporated by reference.\nThe relationships between various cytokines are complex. As will be seen from the references cited herein, as the level of a particular cytokine increases or decreases, this can affect the levels of other molecules produced by a subject, either directly or indirectly. Among the affected molecules are other cytokines.\nThe lymphokine IL-9, previously referred to as \"P40,\" is a T-cell derived molecule which was originally identified as a factor which sustained permanent antigen independent growth of T4 cell lines. See, e.g., Uyttenhove et al., Proc. Natl. Acad. Sci. 85: 6934 (1988), and Van Snick et al., J. Exp. Med. 169: 363 (1989), the disclosures of which are incorporated by reference, as is that of Simpson et al., Eur. J. Biochem. 183: 715 (1989).\nThe activity of IL-9 was at first observed on restricted T4 cell lines, failing to show activity on CTLs or freshly isolated T cells. See, e.g., Uyttenhove et al., supra, and Schmitt et al., Eur. J. Immunol. 19: 2167 (1989). This range of activity was expanded when experiments showed that IL-9 and the molecule referred to as T cell growth Factor III (\"TCGF III\") are identical to MEA (Mast Cell Growth Enhancing Activity), a factor which potentiates the proliferative response of bone marrow derived mast cells to IL-3, as is described by Hultner et al., Eur. J. Immunol. and in U.S. patent application Ser. No. 498,182 filed Mar. 23, 1990, the disclosures of both being incorporated by reference herein. It was also found that the human form of IL-9 stimulates proliferation of megakaryoblastic leukemia. See Yang et al., Blood 74: 1880 (1989). Recent work on IL-9 has shown that it also supports erythroid colony formation (Donahue et al., Blood 75(12): 2271-2275 (Jun. 15, 1990)); promotes the proliferation of myeloid erythroid burst formation (Williams et al., Blood 76: 306-311 (Sep. 1, 1990); and supports clonal maturation of BFU-E's of adult and fetal origin (Holbrook et al., Blood 77(10): 2129-2134 (May 15, 1991)). Expression of IL-9 has also been implicated in Hodgkins's disease and large cell anaplastic lymphoma (Merz et al., Blood 78(8): 1311-1317 (Sep. 1, 1990). Genetic analyses of mice that were susceptible or resistant to the development of bronchial hyperresponsiveness have unraveled a linkage with the IL-9 gene as well as a correlation between IL-9 production and susceptibility in this model (Nicolaides et al., Proc. Natl. Acad. Sci. USA, 94, 13175-13180, 1997). Human genetic studies also point to the IL-9 and IL-9R genes as candidates for asthma (Doull et al., Am. J. Respir. Crit. Care Med., 153, 1280-1284, 1996; Holroyd et al., Genomics 52, 233-235, 1998). Secondly, IL-9 transgenic mice allowed for the demonstration that increased IL-9 expression result in lung mastocytosis, hypereosinophilia, bronchial hyperresponsiveness and high levels of IgE (Temann et al., J. Exp. Med. 188, 1307-1320, 1998; Godfraind et al., J. Immunol. 160, 3989-3996, 1998; McLane et al., Am. J. Resp. Cell. Mol. 19:713-720 (1999). Taken together, these observations strongly suggest that IL-9 plays a major role in this disease Additional work has implicated IL-9 and muteins of this cytokine in asthma and allergies. See, e.g. PCT Application US96/12757 (Levitt, et al), and PCT US97/21992 (Levitt, et al), both of which are incorporated by reference..\nIL-9 is known to affect the levels of other molecules in subjects. See Louahed et al., J. Immunol. 154: 5061-5070 (1995; Demoulin et al., Mol. Cell. Biol. 16: 4710-4716 (1996), both incorporated by reference. It will be recognized that the molecules affected have their own functions in biological systems. For example, Demoulin et al. show that many of the known activities of IL-9 are mediated by activation of STAT transcription factors. As such, there is continued interest in trying to identify molecules whose presence and/or level is affected by other molecules, such as cytokines.\nThe disclosure which follows describes such molecules. It was found that nucleic acid molecules encoding the proteins of the invention were expressed in the presence of IL-9, but not in its absence. Hence, these molecules are, inter alia, \"markers\" for the expression or effect of IL-9 in a subject. The molecules are referred to as T Cell Derived Inducible Factors or \"TIFs\" hereafter. These and other features of the invention will be seen in the disclosure which follows."}
-{"text": "1. Field of the Invention\nThe invention in general relates to the fabrication of layered superlattice materials, and more particularly to a fabrication method using unreactive gas annealing and a low temperature pretreatment to reduce exposure of an integrated circuit to oxygen at elevated temperature.\n2. Statement of the Problem\nFerroelectric compounds possess favorable characteristics for use in nonvolatile integrated circuit memories. See Miller, U.S. Pat. No. 5,046,043. A ferroelectric device, such as a capacitor, is useful as a nonvolatile memory when it possesses desired electronic characteristics, such as high residual polarization, good coercive field, high fatigue resistance, and low leakage current. Layered superlattice material oxides have been studied for use in integrated circuits. U.S. Pat. No. 5,434,102, issued Jul. 18, 1995, to Watanabe et al., and U.S. Pat. No. 5,468,684, issued Nov. 21, 1995, to Yoshimori et al., describe processes for integrating these materials into practical integrated circuits. Layered superlattice materials exhibit characteristics in ferroelectric memories that are orders of magnitude superior to alternative types of ferroelectric materials, such as PZT and PLZT compounds.\nIntegrated circuit devices containing ferroelectric elements with layered superlattice materials are currently being manufactured. A typical ferroelectric memory cell in an integrated circuit contains a semiconductor substrate and a metaloxide semiconductor field-effect transistor (xe2x80x9cMOSFETxe2x80x9d) in electrical contact with a ferroelectric device, usually a ferroelectric capacitor. A ferroelectric memory capacitor typically contains a thin film of ferroelectric metal oxide located between a first, bottom electrode and a second, top electrode, the electrodes typically containing platinum. Layered superlattice materials comprise metal oxides. In conventional fabrication methods, crystallization of the metal oxides to produce desired electronic properties requires heat treatments in oxygen gas at elevated temperatures. The heating steps in the presence of oxygen are typically performed at a temperature in the range of from 800xc2x0 C. to 900xc2x0 C. for 30 minutes to two hours. As a result of the presence of reactive oxygen at elevated temperatures, numerous defects are generated in the single crystal structure of the semiconductor silicon substrate, leading to deterioration in the electronic characteristics of the MOSFET. Good ferroelectric properties have been achieved in the prior art using process heating temperatures at about 700xc2x0 C. to crystallize layered superlattice material. See U.S. Pat. No. 5,508,226, issued Apr. 16, 1996, to Ito et. al. Nevertheless, the annealing and other heating times in the low-temperature methods disclosed in the prior art are in the range of three to six hours, which may be economically unfeasible. More importantly, the long exposure time of several hours in oxygen, even at the somewhat reduced temperature ranges, results in oxygen damage to the semiconductor substrate and other elements of the CMOS circuit.\nAfter completion of the integrated circuit, the presence of oxides may still cause problems because oxygen from the thin film tends to diffuse through the various materials contained in the integrated circuit and combine with atoms in the substrate and in semiconductor layers forming oxides. The resulting oxides interfere with the function of the integrated circuit; for example, they may act as dielectrics in the semiconducting regions, thereby virtually forming capacitors. Diffusion of atoms from the underlying substrate and other circuit layers into the ferroelectric metal oxide is also a problem; for example, silicon from a silicon substrate and from polycrystalline silicon contact layers is known to diffuse into layered superlattice material and degrade its ferroelectric properties. For relatively low-density applications, the ferroelectric memory capacitor is placed on the side of the underlying CMOS circuit, and this may reduce somewhat the problem of undesirable diffusion of atoms between circuit elements. Nevertheless, as the market demand and the technological ability to manufacture high-density circuits increase, the distance between circuit elements decreases, and the problem of molecular and atomic diffusion between elements becomes more acute. To achieve high circuit density by reducing circuit area, the ferroelectric capacitor of a memory cell is placed virtually on top of the switch element, typically a field-effect transistor (hereinafter xe2x80x9cFETxe2x80x9d), and the switch and bottom electrode of the capacitor are electrically connected by a conductive plug. To inhibit undesired diffusion, a barrier layer is located under the ferroelectric oxide, between the capacitor\"\"s bottom electrode and the underlying layers. The maximum processing temperature allowable with current barrier technology is in the range of from 700xc2x0 C. to 750xc2x0 C. At temperatures above this range, the highest-temperature barrier materials begin to degrade and lose their diffusion-barrier properties. On the other hand, the minimum feasible manufacturing process temperatures of layered superlattice materials used in the prior art is about 800xc2x0 C., which is the temperature at which deposited layered superlattice materials are annealed to achieve good crystallization. Lower annealing temperatures require much longer time periods of exposure to oxygen, which can result in damage to the integrated circuit.\nFor the above reasons, therefore, it would be useful to have a method for fabricating layered superlattice materials in ferroelectric integrated circuits that minimizes the time of exposure to oxygen at elevated temperature, as well as reduces the maximum temperature used.\nThe embodiments of the present invention reduce the time of exposure of the integrated circuit to oxygen gas at elevated temperature, and reduce fabrication processing temperatures.\nIn an important embodiment of the invention, a portion of the time period during which an integrated circuit is heated or annealed at elevated temperature is conducted in an oxygen-free unreactive gas. The oxygen-free gas may be any relatively unreactive gas or mixture of unreactive gases, such as nitrogen and the noble gases, in particular, argon and helium. A useful result of embodiments of the invention is that when a fabrication method includes annealing in unreactive gas for a significant portion of the total annealing time at elevated temperature, then the ferroelectric polarizability of layered superlattice material is as high or higher than the polarizability of layered superlattice material annealed for the same total annealing time in oxygen only.\nThe invention provides a method of fabricating a thin film of layered superlattice material comprising: providing a substrate and a precursor containing metal moieties in effective amounts for spontaneously forming a layered superlattice material; applying the precursor to the substrate; annealing the solid thin film in an unreactive gas at a temperature in a range of from 600xc2x0 C. to 800xc2x0 C.; and annealing the solid thin film in an oxygen-containing gas at a temperature in a range of from 600xc2x0 C. to 800xc2x0 C. The invention contemplates that the annealing in an unreactive gas may either precede or come after the annealing in an oxygen-containing gas. The annealing in an unreactive gas may be conducted for a time period in the range of from 30 minutes to 100 hours. The annealing in an oxygen-containing gas is conducted for a time period in the range of from 30 minutes to two hours. In a preferred embodiment, the annealing in an oxygen-containing gas is conducted for a time period not exceeding 60 minutes. Typically, this step of annealing is conducted in substantially pure O2 gas.\nA further embodiment of the invention includes heating the substrate after applying the precursor, forming a solid film from the precursor. In an embodiment, the heating is conducted at a temperature not exceeding 600xc2x0 C. The heating typically comprises a step of drying the precursor coating on the substrate at a temperature not exceeding 300xc2x0 C. Typically, drying is accomplished by baking in an oxygencontaining gas, preferably in O2 gas.\nIn another important embodiment of the invention, the heating comprises a step of pretreating the substrate after the precursor is applied and before annealing. Typically, pretreating is conducted in an oxygen-containing gas. Pretreating preferably comprises rapid thermal processing (xe2x80x9cRTPxe2x80x9d) of the substrate having the precursor coating. In a typical embodiment, the rapid thermal processing is conducted at a temperature in the range of from 300xc2x0 C. to 600xc2x0 C. for a time period in the range of from 1 minute to 15 minutes. Pretreating also may comprise a hot plate baking of the substrate. Typically, hot plate baking is conducted in air at a temperature in the range of from 300xc2x0 C. to 600xc2x0 C. for a time period in the range of from 1 minute to 15 minutes. Pretreating may also comprise a furnace pre-annealing of the substrate. Typically, furnace pre-anneal is conducted at a temperature in the range of from 300xc2x0 C. to 600xc2x0 C. for a time period in the range of from 1 minute to 15 minutes.\nIn another embodiment of the invention, the substrate comprises a first electrode, and the method includes steps of forming a second electrode on the solid thin film, after the step of annealing, to form a capacitor, and subsequently performing a step of post-annealing. In a preferred embodiment, the first electrode and the second electrode contain platinum. The step of post-annealing is conducted at a temperature not exceeding 800xc2x0 C., preferably for a time period in the range of from 30 minutes to two hours. In one embodiment of the invention, the post-annealing is conducted in an oxygen-containing gas, typically in O2 gas. In a preferred embodiment of the invention, the post-annealing is conducted in oxygen-free unreactive gas, typically N2 gas.\nIn a preferred embodiment of the invention, an electrically conductive barrier layer is formed on the substrate prior to applying the precursor coating.\nThe thin film of layered superlattice material typically has a thickness not exceeding 500 nanometers (xe2x80x9cnmxe2x80x9d). In one embodiment, the layered superlattice material comprises strontium, bismuth and tantalum. In another embodiment, the precursor includes metal moieties for forming a layered superlattice material thin film comprising strontium, bismuth, tantalum and niobium.\nIn accordance with the invention, the precursor comprises a solution of metal organic precursor compounds containing metal atoms contained in the desired layered superlattice material. Preferably, the metal organic compounds are metal 2-ethylhexanoates.\nNumerous other features, objects and advantages of the invention will become apparent from the following description when read in conjunction with the accompanying drawings."}
-{"text": "1. Field of the Invention\nThe present invention relates to a method for controlling the operation of a hard disk drive device based upon a detected value of surrounding atmospheric pressure and a hard disk drive device capable of controlling its operation based upon the detected value of the atmospheric pressure.\n2. Description of Related Art\nAt least one magnetic recording disk is contained in a hard disk drive device, and the data is stored on both surfaces of the magnetic recording disk. One read/write head is provided on each recording surface. In one example, four read/write heads are used with two recording disks; each read/write head is mounted on a front end of a head support arm, and rear ends of four head support arms are so connected to each other that all the read/write heads address the same radial cylinder positions on each of the four recording surfaces of the two recording disks. In this manner, all the read/write heads address the same radial positions of the four recording surfaces which define a cylinder, and as such the radial cylinder positions are called cylinder position.\nA prior hard disk drive device used a contact start/stop scheme in which a head/slider assembly is landed on a rest region or a non-recording region located inside the inner most recording cylinder of a magnetic recording disk during a power off period, and takes off from the surface of the rest region when the magnetic recording disk is rotated. To realize a reliable take off, the surface of the disk must be roughened. The reasons for forming the roughness on the recording surface is that if the surface is made smooth, the head/slider assembly sticks to the recording surface due to interatomic forces, thereby preventing take off.\nAs the diameter of the magnetic recording disk becomes smaller, many efforts have been made to increase the recording density and to improve S/N ratio.\nTo increase the recording density and to improve S/N ratio, a flying height of the head/slider assembly, i.e., a space between the head/slider assembly and the recording surface of the disk has been decreased. In the magnetic recording disk, it was difficult to decrease the height of the head/slider assembly over the roughened recording surface, since the head/slider assembly tended to contact to top portion of the roughened surface.\nA hard disk drive device has been developed which uses a load/unload scheme in which a ramp element 2 is mounted at the peripheral of the magnetic recording disk 1, as shown in the FIGS. 1(A) and 1(B). The surface of the ramp element 2 is ramped, and a front end 3 of a head support arm 4 rides on the ramped surface of the part 2, and moves in the direction of an arrow 6 to the rest position of the ramp element 2, as shown in FIG. 1(B), when the head support arm 4 is rotated around a pivot point 5 in the clockwise direction in FIG. 1(A). When the read/write operation is started, the magnetic recording disk 1 is rotated by a motor, not shown, and the head support arm 4 is moved from the rest position in a direction of an arrow 8, and is moved to a flying position above the desired cylinder of the magnetic recording disk 1 to read the data from the cylinder or write the data into the cylinder.\nThe rotational speed, i.e., the revolution per minute (RPM), of the magnetic recording disk 1 is so designed to generate an air bearing with an appropriate pressure which causes the head/slider assembly 7 of the head support arm 4 to fly above the surface of the magnetic recording disk 1.\nIn this manner, the head/slider assembly 7 does not land on the surface of the disk 1, so that the surface of the disk 1 can be polished to remove the roughness of the surface of the disk used in the contact start/stop scheme, and the flying height of the head/slider arm 4 in the read/write mode and an error recovery mode can be reduced in comparison with the flying height of the contact start/stop scheme.\nDescribing the flying height of the head/slider assembly 7 in the read/write mode and the error recovery mode with reference to FIG. 2, in the read/write mode, the flying height of the head/slider assembly 7 is maintained at the height H1 by maintaining the first rotation speed. If the hard disk drive device senses a read error or a write error, the hard disk drive device enters the error recovery mode by reducing the rotational speed of the disk 1 to a second rotational speed which is lower than the first rotational speed to maintain the height H2. This height H2 is so selected to bring the lower surface of the head/slider assembly 7 close to the surface of the disk 1 to wipe off any dusts or residual material on the surface of the magnetic recording disk 1.\nThe problem found by the inventors of the present invention is that the head/slider assembly 7 tends to stick to the surface of the magnetic recording disk 1 when the hard disk drive device enters the error recovery mode while the hard disk drive device is used in an airplane flying at a high altitude."}
-{"text": "It is generally accepted that DNA is the crucial target in the cytotoxic effects of ionising radiation. There is considerable evidence to support the view that DNA double-stranded (ds) breaks are particularly important. The DNA damage results from both direct ionisation in the DNA molecule (direct effect) and by indirect effects mediated by the radiolysis products of water. Carbon-centred radicals on the deoxyribose moiety of DNA are thought to be important precursors of strand breaks. Ionising radiation also induces damage in DNA bases. If the level of cellular DNA damage is sufficient, the consequence of irradiation is cell kill, and thus ionising radiation is used as a mode of cancer therapy. For irradiated normal tissues, the cell killing can result in temporary or permanent impairment of tissue and organ function. The extent of such effects is dependant upon the radiation dose, and if sufficient, can be lethal to the organism. For humans and other animals, hematopoiesis is the most radiosensitive organ/function, followed by the gastrointestinal mucosa. Even if radiation induced DNA damage is sublethal, mutagenic lesions can have serious long term consequences, including carcinogenesis.\nThe medical strategies or countermeasures aimed at reducing the extent of radiation-induced effects are broadly described as radioprotectors (which to be effective, generally need to be administered prior to radiation exposure), mitigants/mitigators (which can be effective if administered after irradiation, but before the appearance of symptoms), and treatments which are generally administered after the appearance of symptoms. A sub-class of the prophylactic radioprotectors are drugs that reduce the extent of the initial radiation-induced DNA damage, and it is this sub-class that is the major focus of the present invention.\nThe commercial potential of radioprotectors resides primarily in two distinct arenas. One of these relates to the need to protect normal tissues in cancer radiotherapy patients, and the other concerns the need to assuage the consequences of unplanned irradiation associated with civil scenarios, such as radiation accidents and radiation terrorism, as well as irradiation in military contexts. This second scenario would also include planned exposure to ionising radiation in medical diagnostic procedures, in which administration of radioprotectors could ameliorate the health risks associated with low or modest radiation doses, without interfering with diagnostic imaging processes.\nThe treatment of tumours with ionising radiation (hereinafter referred to as \u201ccancer radiotherapy\u201d) is used extensively in cancer therapy. The goal of such treatment is the destruction of tumour cells and inhibition of tumour cell growth presumably through DNA damage, while minimising damage to non-tumour cells and tissues. The potential for damage to non-tumour cells in the vicinity of the tumour limits the radiation dose that can be administered, which in turn often limits the effectiveness of radiotherapy against certain tumours. This is especially the case in relation to brain tumours and tumours in the abdominal cavity.\nCancer radiotherapy is a very significant public health activity. Given the incidence of cancer in the population and the international assessment that more than 50% of cancer patients benefit from inclusion of radiotherapy in their treatment, more than 10% of the population are likely to experience cancer radiotherapy in their lifetime.\nThe dominant consideration in prescribing radiation doses for cancer radiotherapy is the assessment of tolerance of the most radiosensitive normal tissues/organs in the treatment field. This assessment, together with the expected radiation dose required to eradicate a tumour, often determines whether the treatment strategy is aimed at cure or palliation. In many cases, the maximum tolerable doses are insufficient to eradicate the tumour. This dilemma is embodied in the concept of therapeutic ratio, which represents the ratio of probabilities of tumour control versus normal tissue morbidity. Approaches to improving the therapeutic ratio include:                (a) optimising the physical targeting of the radiation to the tumour;        (b) fractionation of the radiation dose; and        (c) the use of radiomodifiers (which includes both radioprotectors and radiosensitisers, the latter of which can be used to increase the level of cell kill per unit of radiation dose).        \nImproving the physical delivery of radiation has had a considerable impact on the practice of radiotherapy. For example, increasing the energy of x-ray photons from several hundred kilovolts to the present-day megavoltage beams enables the zone of maximum radiation dose to be set at depths of several centimetres, whereas with the older machines the maximum dose was near the skin surface. There are a number of more sophisticated approaches to \u201ctailoring\u201d treatment beams in various stages of development and implementation. Brachytherapy, the use of implanted radioactive sources rather than external beams, is a further approach to improving the physical dose distribution.\nAlmost without exception, curative external beam radiotherapy involves fractionation of the radiation dose. An example of a conventional schedule would be a total of 60 Gray given in thirty 2 Gray fractions. Since cells have the capacity to repair radiation damage between fractions, the fractionated treatment results in much less cell-kill than a single dose of 60 Gray. However, normal cells generally have a greater repair capacity than do tumour cells, so the \u201csparing\u201d effect of fractionation is more marked for normal tissues. In short, fractionation improves the therapeutic ratio.\nExploration of radiomodifiers such as radioprotectors and radiosensitisers has focussed on hypoxic cell sensitisers such as metranidazole and misonidazole. Radioprotectors have received much less attention than radiosensitisers at the clinical level. The nuclear era spawned considerable effort in the development of radioprotectors with more than 4000 compounds being synthesised and tested at the Walter Reed Army Institute of Research in the United States of America in the 1960's. With the exception of a compound that was called WR2728 (later called Ethyol and now known as Amifostine) none of the compounds have proved useful for cancer radiotherapy, and even WR2728 was considered too toxic for administration in either the military or industrial contexts (i.e., protection against total body irradiation). More recently, for example, Metz and co-workers (Metz et al, Clin Cancer Res. 10, 6411-17, 2004) (15) developed the radioprotective compound known as TEMPOL, which demonstrates only limited efficacy even at very high concentrations, and Burdelya and colleagues (Burdelya et al Science 320, 226-30, 2008) (16) developed the compound known as the TOLL receptor agonist which suffers from the necessity for it to be administered systemically.\nIt is important to note the interplay between the three approaches (a)-(c), above, to improving the therapeutic ratio. A combination of improved physical targeting, fractionation and radiomodifiers could transform the intent in some radiotherapy situations from palliative to curative. For curative schedules, successful application of radiomodifiers would relax the requirement for fractionation and hence reduce overall costs of treatment, which to a large extent is proportional to the number of treatment fractions per patient.\nA particularly important role for radioprotectors has emerged from the recognition that accelerated repopulation of tumour cells during radiotherapy can seriously compromise the effectiveness of treatment. The main consequences of this have been as follows:                (i) The development of accelerated treatment schedules to reduce the overall time of radiotherapy treatment. In such accelerated schedules, acute reactions are a particular problem. For example, acute oral mucositis in head and neck cancer patients indicates a clear need for radioprotectors.        (ii) The recognition that the interruption of radiotherapy treatment due to normal tissue reactions will reduce the probability of tumour control. Accordingly, the use of radioprotectors to prevent toxicity-induced treatment interruption would be clearly beneficial.        \nThe events of 11 Sep. 2001 prompted assessments of vulnerability to many types of terrorism scenarios, amongst which is a collection described as radiological terrorism. An example is the so-called \u201cdirty bomb\u201d involving dispersal of some form a radioactivity with conventional explosive. Whilst attention is focused on the acute radiation syndrome (ARS; also referred to as \u201cradiation sickness\u201d), which describes the consequences of whole-body exposure to radiation doses greater than 1 Gy, there are also concerns about the longer-term effects of low doses, namely radiation-induced mutagenesis and carcinogenesis (1). This general situation, and the realisation that no prophylactic agents are available to provide protection against exposure to ionising radiation, has generated significant research and political activity.\nThe mean lethal dose of radiation required to kill 50% of humans 60 days after whole-body irradiation (LD50/60) is between 3.25 and 4 Gy without supportive care, and 6-7 Gy when antibiotics and transfusion support are provided (1). The mortality is largely attributed to the haematopoietic syndrome, a consequence of hypoplasia or aplasia of the bone marrow. Cytopenias develop as a result of radiation-induced and normal attrition of mature functional cells, combined with the failure of replacement because of radiation-induced depletion of haematopoietic stem cells and progenitors. The time and extent of cytopenia generally correlate with radiation dose and prognosis, but the kinetics of depletion and recovery of blood cells also varies between the erythropoiesis, myelopoiesis and thrombopoiesis lineages, thrombopoiesis being the slowest.\nThe gastrointestinal syndrome results from ablation of stem cells in intestinal crypts, which in turn leads to denudation of the intestinal mucosa. This injury occurs after whole-body doses in the range of 3-15 Gy and in rodents doses at the upper end of this range usually result in death within about 1 week after irradiation.\nCountermeasures against unplanned irradiation include a wide range of potential molecular and cellular interventions. However, the mechanistic simplicity of chemical radioprotection\u2014that is, reduction of radiation-induced DNA damage\u2014is attractive because of its widespread potential. In this context, the possible need for protection of individuals at risk of exposure to low radiation doses, to thereby minimise long-term radiation effects such as mutagenesis and carcinogenesis, is particularly important. Such individuals would include emergency personnel involved in response to unplanned exposures, as well as those subject to occupational exposure to ionising radiation.\nA further group would be patients exposed to ionizing radiation during diagnostic medical procedures conducted in diagnostic radiology and nuclear medicine departments of hospitals and outpatient facilities.\nThe radioprotective properties of the minor groove binding DNA ligand Hoechst 33342 were first described by Smith, P. J. and Anderson, C. O. (2), who used clonogenic survival assays of irradiated cultured cells. Young, S. D. and Hill, R. P. (3) reported similar effects in cultured cells, but extended their studies to in vivo experiments. They concluded that the lack of radioprotection in their in vivo experiments was due to insufficient levels of Hoechst 33342 being delivered to target cells following intravenous injection. The findings of Hill and Young underline an important requirement for effective radioprotectors, namely potency. If the radioprotector is more potent, then it is more likely to achieve the required concentrations in an in vivo setting.\nThere is another aspect to be considered apart from potency. The concentration required for radioprotection must be non-toxic regardless of the potency of the radioprotector. If the radioprotector is delivered systemically, then this lack of toxicity requirement includes not just the cells and tissues to be protected from the radiation, but extends to the toxicity to the subject as a whole. In the case of Hoechst 33342 toxicity limits the extent to which it is useful as a radioprotector.\nThere is also a substantial conceptual problem in using radioprotectors in cancer radiotherapy. In attempting to decrease the effect of radiation on normal tissues by application of radioprotectors, there is a fear that some of the radioprotector will reach the tumour, thereby compromising tumour cell kill. The existing radioprotectors, e.g. WR2721 (Amifostine) and its active metabolite WR1065, are relatively small, diffusible molecules which do not avidly bind to tissue components and can therefore penetrate effectively through cell layers, so that they can reach the tumour via the circulation.\nThere is a need for radioprotectors that have limited penetration through cell layers. Such a property enables radioprotectors to be applied locally or topically to critical radiosensitive normal tissues in the vicinity of the tumour. Limited penetration restricts the extent to which the radioprotector reaches the capillary bed and is taken up into the circulation thereby reaching the tumour by systemic delivery in sufficient concentrations to confer significant radioprotection to the tumour.\nThe limited diffusion of DNA-binding ligands such as Hoechst 33342 through cell layers is known and has been exploited in mapping the location of cells in multi-cellular spheroids and in vivo, with respect to perfusion. Thus perfusion of Hoechst 33342 is considered a surrogate marker for perfusion of oxygen. In addition to restricting access to the tumour by systemic uptake following local or topical application to normal tissues, there is a further potential advantage of limited penetration in the context of cancer radiotherapy. This advantage stems from the view that the vasculature, in particular the endothelial cells, are the critical targets that determine the damaging effects of radiation. Furthermore, most radioresistant cells in the tumour are those viable cells that are most distant from the capillaries. The radioresistance of these cells is due to their hypoxic state, which in turn reflects their remoteness from the capillaries.\nConsequently, radioprotectors having limited diffusion, when administered intravenously, will be delivered more efficiently to critical radiosensitive cells in normal tissues, than to the hypoxic subpopulation of cells in tumours which limit the effectiveness of radiotherapy generally. Thus, the use of such radioprotectors would be expected to enable higher radiation doses to be used, with increased probability of killing the hypoxic cells in the tumour.\nHowever, the potential of the combination of these radiobiological features and the characteristics of DNA-binding radioprotectors can only be useful in cancer radiotherapy provided that an over-riding and necessary requirement of the radioprotectors exists, namely that the radioprotectors are sufficiently potent as to confer demonstrable radioprotection at non-toxic concentrations, when applied topically or systemically. A further practical requirement is that the extent of the limited penetration is sufficient to prevent significant systemic uptake following topical application, but not so pronounced so as to prevent sufficient concentrations from reaching the cells that determine the radiosensitivity of the tissue to be protected from the effects of ionising radiation, by topical or local application.\nThe extent of radioprotection (in the contexts of both cancer radiotherapy and protection from unplanned radiation exposure) is generally described in terms of dose modification factor (DMF), which is defined as the ratio of radiation doses required to produce the equivalent radiation-induced effect (molecular, cellular or in vivo endpoint) in the presence and absence of the radioprotector. When the radioprotective effect is observed on the basis of an in vivo endpoint, mechanisms other than modification of the initial radiation-induced damage may be involved. For example, for both the haematopoietic syndrome and the gastrointestinal syndrome, infection plays an important role in ultimate mortality, as a consequence of neutropenia and breach of the intestinal mucosal barrier, respectively. Thus, some immunostimulants have potential as mitigators of the radiation response. Immunostimulants can also be effective post-irradiation.\nInternational patent publication No. WO97/04776 and the subsequent publication by Martin et at (4) disclose certain bibenzimidazole compounds characterised by substitution with sterically hindering and electron donating groups. Although these compounds demonstrate strong radioprotective activity there is scope to reduce the inherent cytotoxicity of compounds of this general class. The challenge, however, is to do so while retaining, and preferably improving, radioprotective activity (measured as dose modification factor). The disclosures of WO97/04776 are included herein in their entirety by way of reference.\nInternational patent publication No. WO/2008/074091 also discloses bibenzimidazole compounds substituted with fluorine and/or chlorine and which, relative to known radioprotector compounds such as those of International patent publication No. WO97/04776, exhibit reduced cytotoxicity activity. While the cytotoxicity of the fluorine and chlorine substituted bibenzimidazole compounds was improved there is still a need for development of alternative radioprotective compounds, and preferably compounds that can be used in cancer radiotherapy, in protection of biological material from effects of radiation exposure and/or in protection of humans or animals from the effects of unplanned irradiation, which demonstrate low cytotoxicity but that retain radioprotective potency, and preferably that penetrate through cell layers to a limited extent. In particular it is desirable in some contexts that such compounds can be administered topically to protect tissues such as the skin, oral mucosa, oesophageal mucosa, rectal mucosa, vaginal mucosa and bladder epithelium, as well as parenterally to protect organs such as the lung and brain."}
-{"text": "Uroguanylin, guanylin and bacterial ST peptides are structurally related peptides that bind to a guanylate cyclase receptor and stimulate intracellular production of cyclic guanosine monophosphate (cGMP) (1-6). This results in the activation of the cystic fibrosis transmembrane conductance regulator (CFTR), an apical membrane channel for efflux of chloride from enterocytes lining the intestinal tract (1-6). Activation of CFTR and the subsequent enhancement of transepithelial secretion of chloride leads to stimulation of sodium and water secretion into the intestinal lumen. Therefore, by serving as paracrine regulators of CFTR activity, cGMP receptor agonists regulate fluid and electrolyte transport in the GI tract (1-6; U.S. Pat. No. 5,489,670).\nThe process of epithelial renewal involves the proliferation, migration, differentiation, senescence, and eventual loss of GI cells in the lumen (7,8). The GI mucosa can be divided into three distinct zones based on the proliferation index of epithelial cells. One of these zones, the proliferative zone, consists of undifferentiated stem cells responsible for providing a constant source of new cells. The stem cells migrate upward toward the lumen to which they are extruded. As they migrate, the cells lose their capacity to divide and become differentiated for carrying out specialized functions of the GI mucosa (9). Renewal of GI mucosa is very rapid with complete turnover occurring within a 24-48 hour period (9). During this process mutated and unwanted cells are replenished with new cells. Hence, homeostasis of the GI mucosa is regulated by continual maintenance of the balance between proliferation and apoptotic rates (8).\nThe rates of cell proliferation and apoptosis in the gut epithelium can be increased or decreased in a wide variety of different circumstances, e.g., in response to physiological stimuli such as aging, inflammatory signals, hormones, peptides, growth factors, chemicals and dietary habits. In addition, an enhanced proliferation rate is frequently associated with a reduction in turnover time and an expansion of the proliferative zone (10). The proliferation index has been observed to be much higher in pathological cases of ulcerative colitis and other GI disorders (11). Thus, intestinal hyperplasia is the major promoter of gastrointestinal inflammation and carcinogenesis.\nIn addition to a role for uroguanylin and guanylin as modulators of intestinal fluid and ion secretion, these peptides may also be involved in the continual renewal of GI mucosa. Previously published data in WO 01/25266 suggests a peptide with the active domain of uroguanylin may function as an inhibitor of polyp development in the colon and may constitute a treatment of colon cancer. However, the mechanism by which this is claimed to occur is questionable in that WO 01/25266 teaches uroguanylin agonist peptides that bind specifically to a guanylate cyclase receptor, termed GC-C, that was first described as the receptor for E. coli heat-stable enterotoxin (ST) (4). Knockout mice lacking this guanylate cyclase receptor show resistance to ST in intestine, but effects of uroguanylin and ST are not disturbed in the kidney in vivo (3). These results were further supported by the fact that membrane depolarization induced by guanylin was blocked by genistein, a tyrosine kinase inhibitor, whereas hyperpolarization induced by uroguanylin was not effected (12,13). Taken together these data suggest that uroguanylin also binds to a currently unknown receptor, which is distinct from GC-C.\nOther papers have reported that production of uroguanylin and guanylin is dramatically decreased in pre-cancerous colon polyps and tumor tissues (14-17). In addition, genes for both uroguanylin and guanylin have been shown to be localized to regions of the genome frequently associated with loss of heterozygosity in human colon carcinoma (18-20). Taken together, these findings indicate that uroguanylin, guanylin and other peptides with similar activity may be used in the prevention or treatment of abnormal colon growths. This proposal is bolstered by a recent study demonstrating oral administration of uroguanylin inhibits polyp formation in mice (15,16).\nUroguanylin and guanylin peptides also appear to promote apoptosis by controlling cellular ion flux. Alterations in apoptosis have been associated with tumor progression to the metastatic phenotype. While a primary gastrointestinal (GI) cancer is limited to the small intestine, colon, and rectum, it may metastasize and spread to such localities as bone, lymph nodes, liver, lung, peritoneum, ovaries, brain. By enhancing the efflux of K+ and influx of Ca++, uroguanylin and related peptides may promote the death of transformed cells and thereby inhibit metastasis.\nOne of the clinical manifestations of reduced CFTR activity is the inflammation of airway passages (21). This effect may be due to CTFR regulating the expression of NF-KB, chemokines and cytokines (22-25). Recent reports have also suggested that the CFTR channel is involved in the transport and maintenance of reduced glutathione, an antioxidant that plays an important role in protecting against inflammation caused by oxidative stress (39). Enhancement of intracellular levels of cGMP by way of guanylate cyclase activation or by way of inhibition of cGMP-specific phosphodiesterase would be expected to down-regulate these inflammatory stimuli. Thus, uroguanylin-type agonists should be useful in the prevention and treatment of inflammatory diseases of the lung (e.g., asthma), bowel (e.g., ulcerative colitis and Crohn's disease), pancreas and other organs.\nOverall, it may be concluded that agonists of guanylate cyclase receptor such as uroguanylin have potential therapeutic value in the treatment of a wide variety of inflammatory conditions, cancer (particularly colon cancer) and as anti-metastatic agents. The development of new agonists is therefore of substantial clinical importance."}
-{"text": "Duplex stainless steel refers to ferritic austenitic steel alloy. Such steels have a microstructure comprising ferritic and austenitic phases. The duplex steel alloy, to which the invention pertains, is characterized by a high content of Cr and N and a low content of Ni. Background references in this respect include WO 95/00674 and U.S. Pat. No. 7,347,903. The duplex steels described therein are highly corrosion resistant and can therefore be used, e.g., in the highly corrosive environment of a urea manufacturing plant.\nUrea (NH2CONH2) can be produced from ammonia and carbon dioxide at elevated temperature (typically between 150\u00b0 C. and 250\u00b0 C.) and pressure (typically between 12 and 40 MPa) in the urea synthesis section of a urea plant. In this synthesis, two consecutive reaction steps can be considered to take place. In the first step, ammonium carbamate is formed, and in the next step, this ammonium carbamate is dehydrated so as to provide urea, The first step (i) is exothermic, and the second step can be represented as an endothermic equilibrium reaction (ii):2NH3+CO2\u2192H2N\u2014CO\u2014ONH4\u2003\u2003(i)H2N\u2014CO\u2014ONH4H2N\u2014CO\u2014NH2+H2O\u2003\u2003(ii)\nIn a typical urea production plant, the foregoing reactions are conducted in a urea synthesis section so as to result in an aqueous solution comprising urea. In one or more subsequent concentration sections, this solution is concentrated to eventually yield urea in the form of a melt rather than a solution. This melt is further subjected to one or more finishing steps, such as prilling, granulation, pelletizing or compacting.\nA frequently used process for the preparation of urea according to a stripping process is the carbon dioxide stripping process, as for example described in Ullmann's Encyclopedia of Industrial Chemistry, Vol. A27, 1996, pp 333-350. In this process, the synthesis section is followed by one or more recovery sections. The synthesis section comprises a reactor, a stripper, a condenser and, preferably but not necessarily, a scrubber in which the operating pressure is in between 12 and 18 MPa, such as in between 13 and 16 MPa. In the synthesis section, the urea solution leaving the urea reactor is fed to a stripper in which a large amount of non-converted ammonia and carbon dioxide is separated from the aqueous urea solution.\nSuch a stripper can be a shell- and tube-heat exchanger in which the urea solution is fed to the top part at the tube side and a carbon dioxide feed, for use in urea synthesis, is added to the bottom part of the stripper. At the shell side, steam is added to heat the solution. The urea solution leaves the heat exchanger at the bottom part, while the vapor phase leaves the stripper at the top part. The vapor leaving said stripper contains ammonia, carbon dioxide, inert gases and a small amount of water.\nSaid vapor is condensed in a falling film type heat exchanger or a submerged type of condenser that can be a horizontal type or a vertical type. A horizontal type submerged heat exchanger is described in Ullmann's Encyclopedia of Industrial Chemistry, Vol. A27, 1996, pp 333-350. The formed solution, which contains condensed ammonia, carbon dioxide, water and urea, is recirculated together with the non-condensed ammonia, carbon dioxide and inert vapor.\nThe processing conditions are highly corrosive, particularly due to the hot carbamate solution. In the past, this presented a problem in the sense that the urea manufacturing equipment, even though made from stainless steel, would corrode and be prone to early replacement.\nThis has been resolved, particularly by making the equipment, i.e. the relevant parts thereof subjected to the mentioned corrosive conditions, from a duplex steel described in WO 95/00674 (also known by the trademark of Safurex\u00ae). However, even though the foregoing reflects a major advancement in urea production, a particular problem exists in the stripper. A typical carbamate stripper comprises a plurality (several thousand) of tubes. Through the tubes, a liquid film runs downwards whilst stripping gas (typically CO2) runs upwards. Provisions are generally made to ensure that all tubes have the same load of liquid so as to have a flow of the liquid at the same speed. For, if the liquid does not flow through all of the tubes at the same speed, the efficiency of the stripper is reduced. These provisions comprise a liquid distributor, generally in the form of a cylinder with small holes in it.\nIt has been experienced that the liquid distributors need a relatively frequent replacement. Particularly, the size and shape of the holes changes with time, apparently as a result of corrosion, despite the fact that the liquid distributors are made from corrosion-resistant duplex steel as mentioned above. Thus, the affected distributors result in a different throughput of liquid in the stripper, as a result of which the desired equal loading of the stripper's tubes is less efficient.\nIt is therefore desired in the art to provide a corrosion resistant material that would provide the liquid distributors in the stripper with a better corrosion endurance."}
-{"text": "The present invention relates generally to the packaging of integrated circuits. More particularly, the invention relates to dice having integrated pedestals.\nThere are a number of conventional processes for packaging integrated circuits. In many situations it is desirable to incorporate multiple integrated circuit dice into the same package in what is commonly referred to as a multi-chip package. Some multi-chip packages are arranged to stack two or more dice on top of each other. These stacked die packages have several potential advantages including the possibility of a reduced die or package footprint and certain performance advantages. For example, reducing the path length of electrical connections between integrated circuits potentially increases speed and reduces inductance of inter-chip communications.\nOne existing stacked die configuration is generally illustrated in FIG. 1. In this configuration a first die 104 is mounted on a carrier such as a lead frame 106 or a planar substrate. A second die 105 is then adhesively secured to the top surface of the first die 104 thereby creating a stacked die configuration. Bonding wires 108 are used to electrically connect both of the dice 104, 105 to the carrier (lead frame 106) using conventional wire bonding. An encapsulant material 115 is molded over the stacked dice to create a stacked die package 120. This stacked die approach allows the top die 105 to be electrically connected to both the underlying base die 104 and to a substrate or lead frame 106. This arrangement also generally requires that the top die 105 must be smaller than the base die 104 so that the top die 105 does not cover any of the bond pads on the base die 104.\nAnother stacked die approach is generally illustrated in FIG. 2. In this approach, a spacer 130 is adhesively secured to (or dispensed onto) the top surface of the bottom die 104. The top die 105 is then adhesively adhered to the spacer 130. The spacer 130 is sized so that it extends higher than the loop height of the bonding wires 108 used to electrically connect the base die, thereby providing clearance for the bonding wires. The spacer 130 can be formed from virtually any material that is compatible with the dice. By way of example, the spacer may be formed from a variety of materials such as silicon, various ceramics, etc. Alternatively, the spacer may be formed in situ by dispensing an adhesive material having ball like spherical support structures therein on the top surface of the lower die 104 in a region that is interior to the bond pads. The spacer 130 effectively forms a plateau on the bottom die that provides structural support for the top die. With this arrangement, the top die 105 can be virtually any size including the same size or even larger than the base die.\nAlthough the foregoing techniques work well in many applications, in the semiconductor industry, there are continuing efforts to provide more efficient packaging approaches. The packaging arrangements described below are particularly useful in forming stacked multi-chip packages."}
-{"text": "Conventional gold-based catalysts, and particularly nanoporous gold-based catalysts either consist of the nanoporous gold material, or are formed by depositing gold on and/or within metal oxides to form an active structure. While these conventional structures exhibit beneficial activity at relatively low temperatures (e.g. about 0 centigrade), as temperature increases the activity degrades due to the nanoporous gold material forming aggregates (e.g. by sintering), reducing the surface area of the overall structure and thus the catalytic activity thereof.\nThis is particularly undesirable for use in applications to which such catalysts have been earnestly researched. In particular, for an engine (in particular an internal-combustion engine) that operates at a wide range of temperatures, a catalyst that could efficiently treat exhaust gas both in a cold-start phase and at a peak-operation phase when the engine is hot (e.g. several hundred centigrade) would be beneficial."}
-{"text": "The present invention relates to the microwave generation, amplification, and processing arts. It particularly relates to traveling wave tubes for microwave amplifiers and microwave oscillators, and will be described with particular reference thereto. However, the invention will also find application in other devices that operate at microwave frequencies, and in other devices that employ slow wave circuits.\nTraveling wave tubes typically include a slow wave circuit defined by a generally hollow vacuum-tight barrel with optional additional microwave circuitry disposed inside the barrel. An electron source and suitable steering magnets or electric fields are arranged around the slow wave circuit to pass an electron beam through the generally hollow beam tunnel. The electrons interact with the slow wave circuit, and energy of the electron beam is transferred into microwaves that are guided by the slow wave circuit. Such traveling wave tubes provide microwave generation and microwave amplification.\nHeretofore, commercially produced traveling wave tubes have been limited to about 65 GHz. However, future applications call for traveling wave tubes that operate at frequencies of 100 GHz or higher. For space-based applications these high frequency devices will probably be driven at operating voltages of 20 kV or less in accord with presently available space-based electrical power sources.\nConstruction of high frequency traveling wave tubes is difficult using existing traveling wave tube manufacturing techniques. Designs for high frequencies call for microwave circuitry with very small features (for example, repetition periods of less than 0.2 mm), and greatly reduced quantities of dielectric insulation material in the tube to reduce dielectric loading. Moreover, adequate heat sinking becomes an increasingly significant issue as the operating frequency increases.\nIn one known method for manufacturing traveling wave tubes, a fragile three-dimensional microwave circuit shell, such as a metallic helix, is compressively secured in a generally hollow cylindrical barrel. Dielectric rods arranged parallel to the helical axis of the microwave circuit act as standoff insulators that align and secure the compressed microwave circuit shell inside the barrel.\nTo ensure good thermal contact between the components, the compressive forces are substantial. The fragile microwave circuit shell and dielectric rods are compressively secured in the barrel by briefly heating the barrel during insertion to induce temporary thermal expansion of the barrel. The microwave circuit shell/dielectric rods combination has close tolerances with respect to the barrel, and so when the barrel contracts upon cooling the interior components are compressively secured in the barrel. However, the heating and compression can damage the slow wave circuit, and mass production by this method is difficult. Moreover, this technique is not well suited for the small structures used in devices that are preferred for 100 GHz or higher operation.\nTo achieve features on the delicate scale called for in high frequency operation, lithographic techniques are regularly used in the electronics industry. However, these techniques are generally applied to planar wafer substrates of silicon or other semiconductor materials. Lithographic techniques are not readily adapted to produce the types of finely detailed three-dimensional structures called for in traveling wave tubes designed for high frequency operation.\nTo reduce dielectric loading, the dielectric rods can be replaced by thin dielectric standoff chips of natural diamond. In one constructed device described in A Millimeter-Wave TunneLadder TWT (D. Wilson, NASA Contract Report 182183, 1988), diamond chips with heights of 250 microns each were used in a traveling wave tube that operated at 28 GHz. However, this device has not been replicated to date due to the cost of natural diamond and the assembly difficulties, especially relating to positioning of the diamond chips. Moreover, scaling such a device up to 100 GHz or higher frequency would call for a large number of diamond chips (e.g., about 80-150 diamond chips for a 2-3 cm long traveling wave tube designed for 100 GHz operation) each having a height of about 75 microns or thinner. These reduced dimensions and increased numbers of diamond chips would further exacerbate an already difficult manufacturing process.\nThe present invention provides an improved apparatus and method."}
-{"text": "The present invention relates to telecommunication systems and more particularly to digital hierarchy systems (SDH). In such systems data is switched within an individual station to the appropriate output port. Also alternate routes are provided between stations by which data can be transmitted from an originating stations port to a receiver stations port so that failure of one route, e.g. due to accidental damage, will not prevent the data from reaching the receiver stations port via the other route.\nThe data is sent in the form of two packages or so-called virtual containers which essentially are made up of two sections, one being the actual data which is transmitted, i.e. the soiled payload, and the other being data which is concerned with the integrity of the data, referred to hereinafter as overhead or ancillary data.\nOne of the requirements for a system of the kind outlined above is that there must be some means for checking whether the package is correctly switched (or routed) at each individual station along the overall path of the data. Essentially this is achieved by use of the ancillary or overhead data area at each station in order to include an identity, and other monitoring information, to ensure correct switching.\nIn our co-pending application no. 9301575 we disclose an arrangement where the package is identified in terms of the channel number of the data, and the receiving stations card number. However there is a potential problem with such an arrangement in situations where the package travels by the alternative route referred to earlier. In the event of failure of the first route the package will acquire a new channel and card number. This means that the destination port of the station will not recognise the package because it will be looking for a package which has an identification which is characteristic of the first route. This problem is a direct result of the use of the alternate route protection identified earlier.\nThe present invention is concerned with overcoming this problem."}
-{"text": "In general, heat is an enemy for the IC chip and it is necessary that the internal temperature thereof does not exceed the maximum allowable junction temperature. The electric power consumption per operation area is large in the semiconductor device such as a power transistor or a semiconductor rectifier element. Therefore, the generated heat amount cannot be sufficiently released with only the heat amount released from a case (package) and a lead of the semiconductor device. It is feared that the internal temperature of the device may be raised to cause thermal destruction.\nThis phenomenon also occurs in the same manner in the IC chip which carries a CPU. The amount of heat generation is increased during the operation in proportion to the improvement in clock frequency. It is an important matter to make the thermal design in consideration of the heat release.\nIn the thermal design for preventing the thermal destruction or the like, element design or package design is made on condition that a heat sink having a large heat release area is secured to a case (package) of the IC chip.\nIn general, a metal material such as copper and aluminum, which has a good thermal conductivity, is used as a material for the heat sink.\nRecently, the IC chip such as CPU and memory is in a trend that the IC chip itself has a large size in accordance 10, with the high degree of integration of the element and the enlargement of the element-forming area, while it is intended to drive the IC chip at low electric power for the purpose of low electric power consumption. When the IC chip has such a large size, it is feared that the stress caused by the difference in thermal expansion between the semiconductor substrate (silicon substrate or GaAs substrate) and the heat sink is increased, and that the peeling-off phenomenon and the mechanical destruction occur in the IC chip.\nIn order to avoid such an inconvenience, for example, it may be pointed out that the low electric power driving of the IC chip should be realized, and the heat sink material should be improved. The low electric power driving of the IC chip is realized in the level of not more than 3.3 V at present and the TTL level (5 V) which has been hitherto used as the power source voltage becomes obsolete.\nAs for the constitutive material for the heat sink, it is insufficient to consider only the thermal conductivity. It is necessary to select a material which has a coefficient of thermal expansion approximately identical with those of silicon and GaAs, which are used as the semiconductor substrate, while having a high thermal conductivity at the same time.\nA variety of reports have been made in relation to the improvement of the heat sink material, including, for example, a case in which aluminum nitride (AlN) is used and a case Cu (copper)-W (tungsten) is used. AlN is excellent in balance between the thermal conductivity and the thermal expansion. Especially, the coefficient of thermal expansion of AlN is approximately coincident with the coefficient of thermal expansion of Si. Therefore, AlN is preferred as a heat sink material for a semiconductor device in which a silicon substrate is used as the semiconductor substrate.\nCu\u2014W is a composite material having both of the low thermal expansion of W and the high thermal conductivity of Cu. Further, Cu\u2014W is mechanically machined with ease. Therefore, Cu\u2014W is preferred as a constitutive material for a heat sink having a complicated shape.\nOther instances have been suggested, wherein metal Cu is contained in a ratio of 20 to 40% by volume in a ceramic base material containing a major component of SiC (conventional technique 1, see Japanese Laid-Open Patent Publication No. 8-279569), and wherein a powder-sintered porous member of an inorganic substance is infiltrated with Cu by 5 to 30% by weight (conventional technique 2, see Japanese Laid-Open Patent Publication No. 59-228742).\nThe heat sink material concerning the conventional technique 1 is produced in the powder formation in which a green compact of SiC and metal Cu is formed to produce a heat sink. Therefore, the coefficient of thermal expansion and the coefficient of thermal conductivity represent only theoretical values. It is impossible to obtain the balance between the coefficient of thermal expansion and the coefficient of thermal conductivity required for actual electronic parts etc.\nThe conventional technique 2 uses a low ratio of Cu with which the powder-sintered porous member composed of the inorganic substance is infiltrated. It is feared that a limit may arise to enhance the thermal conductivity thereby.\nOn the other hand, a composite material, which is obtained by combining carbon and metal, has been developed and practically used. However, such a composite material is used, for example, as an electrode for the discharge machining when the metal is Cu. When the metal is Pb, such a composite material is used, for example, as a bearing material. No case is known for the application as a heat sink material.\nThat is, in the present circumstances, the coefficient of thermal conductivity is at most 140 W/mK for the composite material obtained by combining carbon and metal, which cannot satisfy the value of not less than 160 W/mK required for the heat sink material for the IC chip."}
-{"text": "Five unique human hepatitis viruses have been identified (1-5). The hepatitis A virus and hepatitis E virus are enterically transmitted RNA viruses that do not cause chronic liver disease. In contrast, the hepatitis B virus, hepatitis C virus and hepatitis D virus (HBV, HCV and HDV, respectively) are parenterally transmitted and cause chronic infection. They are dangerous contaminants of the blood supply. Recently tests have become readily available for testing for HBV in blood, allowing for the screening for this pathogen and the elimination of infected samples from the blood supply (6).\nConcomitant with the availability of the HBV test came an increase in the proportion of cases of post-transfusion hepatitis due to non-A, non-B (NANB) agents. Until recently, there was no test available for the detection of the NANB agents. The principal NANB agent, HCV, was recently identified by molecular cloning of segments of the HCV genome (3). HCV is an RNA virus related to human flaviviruses and animal pestiviruses (7,8).\nProspective studies of selected counties in the United States by the Centers for Disease Control (CDC) indicate that approximately 170,000 new cases of NANB/HCV infection occur yearly (9). At least 50% of these infections appear to progress to chronic liver disease. Severe sequelae include the development of decompensated cirrhosis necessitating liver transplantation, and development of hepatocellular carcinoma (10,11).\nThe positive-stranded RNA genome of the HCV contains approximately 10,000 nucleotides. The HCV genome acts as a long open reading frame (ORF) capable of encoding a 3,010 amino acid polyprotein precursor from which individual viral proteins, both structural and nonstructural, are produced (7,12-14). There are at least 324 nucleotides at the 5'-end of the ORF which have not yet been shown to encode for protein. Thus, this sequence is referred to as the 5'-non-coding region (7,12-17). Several research groups have reported the nucleotide sequence of either the whole HCV genome or specific subgenomic regions (7,12-22). Comparison of these sequences demonstrates variations in the structural and nonstructural regions (ranging from 9-26%) among different HCV strains. In contrast, the sequences of the 5'-non-coding region appear to have a homology of approximately 99% among different strains (16,17). The 5'-non-coding region also has substantial homology (45-49%) with the equivalent region of animal pestiviruses (7).\nTwo major techniques are currently used to detect HCV infection. The first technique detects antibody produced in response to HCV infection (anti-HCV) (23-28). Since multiple weeks are required for infected patients to develop detectable IgG antibody against HCV antigens, this test is useless in the detection of acute HCV infection. Moreover, studies indicate that antibody testing is associated with both false positive and false negative results (29).\nThese shortcomings in the original assays have spurred development of newer supplemental antibody tests for the diagnosis of HCV infection (30). Preliminary results with supplemental assays indicate a decrease in the frequency of false-positive and negative results. However, false-positive and -negative results still occur and supplemental tests remain unsuitable for detection of acute infection (31).\nThe second technique, detection of HCV RNA by an RNA polymerase chain reaction (PCR), has been limited to research use. The HCV PCR evaluates infection by detecting HCV RNA in blood or tissue extracts through reverse transcription and cDNA amplification (7,32-41). HCV PCR represents a sensitive, direct technique but requires meticulous care (7) to prevent false positive and negative results. The HCV PCR technique, in contrast to antibody tests, can detect circulating HCV RNA during acute infection.\nThe original HCV PCR tests used primers specific for sequences in the non-structural region of the HCV genome (32-36). Subsequently, HCV PCR has been performed using several primers for the 5'-non-coding region in the genome (37,39). In our laboratory we have established HCV PCR for both the nonstructural and 5'-non- coding regions. Our comparative results indicate that the HCV PCR from the 5'-non-coding region is more sensitive in detecting HCV infection (41).\nDespite the success of HCV PCR, the technique has many inherent limitations. First, it is time consuming, expensive and dependent upon meticulous technique. The exquisite sensitivity of PCR makes false positive results due to contamination with exogenous HCV RNA a constant concern (42). Moreover, the variation in both the reverse transcription of HCV RNA to CDNA and the amplification of cDNA make the HCV PCR difficult to quantitate (38,40,42). Recent attempts to overcome these obstacles have resulted in, at best, semi-quantitative assays (38). More importantly, in our experience the efficiency of HCV PCR depends in large part on the specific primers employed. Not only have standards for primers not been developed, but polymerases employed in PCR have different efficiencies. Thus, it will likely be difficult to compare PCR results among different laboratories.\nTo overcome the limitations of current antibody and HCV PCR techniques for detection of HCV infection, it is desirable to develop a test which is highly sensitive, specific, affordable and applicable to the testing of large populations of patients or blood donors. The optimal test would be capable of detecting both acute and chronic infection. Moreover, it would be quantitative to provide information regarding both natural history and the efficacy of current or future antiviral therapies. It would be capable of uniform results. These prerequisites can be fulfilled by a technique to directly detect HCV RNA."}
-{"text": "There are many methods for reproducing stereo-pictures (three dimensional pictures) such that the observer has the full spatial impressions of the object. This is done, e.g., by arranging the two pictures recorded from two different directions, as the two eyes of an observer view from the corresponding directions, on a screen and viewing the pictures through glasses of different color for non-color reproduction. If a color image is required, the different pictures are illuminated with linear polarized light of perpendicular orientation and glasses are used with the appropriate polarization such that the eyes see the objects separately in stereo. One difficulty with the foregoing is that the eyes are focused in an unnatural way, i.e. looking in different directions on alternate frames, if the reproduced pictures have to be next to each other.\nThe ideal method requires that the two reproduced pictures of different polarization be located within the same frame. This can be realized, e.g., by holographic production of the color TV picture. This method, however, is basically limited by the coherence length of the radiation such that the stereo effects can be seen only within a certain depth, e.g., 30 centimeters to 1 meter with the current generation of holographic reproduction equipment. This depth, even if extended, is limited by basic principles of physics. A further difficulty is that the granulation of the transmitted picture requires extremely high resolution and very small grain photographic material must be used. The electronic transmission of such pictures also requires an extraordinarily large bandwidth and minimal signal distortion."}
-{"text": "1. Field of the Invention\nThis invention relates generally to biohazard surveillance systems and more particularly to an adaptive distributed system for the collection and sampling of hazardous particulates.\n2. Description of the Related Art\nThe challenges we face from biological threat agents are increasing. While microbes continue to evolve and biotechnology becomes more powerful, the inherent hazards to humans, plants, and animals from infectious microorganisms are greatly increased by their intentional use by terrorists. The need for faster and better capabilities for warning, response, and cleanup was painfully evident in the case of a small-scale deployment of a noncontagious, naturally occurring anthrax pathogen. Terrorist use of other biological agents may result in far greater loss of life; agents that might be contagious or perhaps engineered for increased virulence and resistance to medical treatment. As microbes evolve and compete for survival, naturally emerging threats must also be quickly identified and distinguished from suspected terrorism. While the focus on bioterrorism is driven primarily by concerns about attacks on humans, attacks on livestock and/or crops can be just as devastating. A recent outbreak of foot-and-mouth disease in Great Britain demonstrates the devastating effect microbes can have on livestock and the consequent effect on food supply and economies. Rogue states have actively explored both animal and plant pathogens as weapons.\nLessons learned from the Persian Gulf War highlighted the need for biological warfare agent detectors and the subsequent solutions improved capability on the battlefield. However, other biological hazard (\u201cbiohazard\u201d) surveillance deficiencies were soon recognized in the aftermath of conflict. \u201cGulf War Syndrome\u201d and other ailments suffered by military personnel revealed a need for compact diagnostic tools with integrated sample-processing and detection capabilities to quickly identify disease-causing agents on and off the battlefield. In 1998, a consolidated approach was begun (at the Army Medical Institute for Infectious Diseases) to develop medical diagnostic systems using a common platform for biohazard identification entitled \u201cThe Common Diagnostic Systems for Biological Threats and Endemic Infectious Diseases.\u201d Research encompassed development of rapid sample-processing methods, identification technologies, reagents and size reduction of laboratory analysis platforms.\nIn 2002, the Department of Defense (DoD) defined a new approach to a common medical test platform for identifying biological warfare agents and pathogens of operational concern. The Joint Biological Agent Identification and Diagnostic System (JBAIDS) exemplifies this approach. JBAIDS will be configured to support reliable, fast, and specific identification of biological agents from a variety of clinical specimens and environmental sources. JBAIDS will enhance force protection by providing commanders with information to determine actions to protect against and avoid contamination and to restore operations following an attack. JBAIDS information will aid medical personnel in determining appropriate treatment, effective preventive medical measures, and medical prophylaxis in response to the presence of biological agents. Required to combat the threat of biological attack faced by U.S. forces deployed worldwide, JBAIDS will also improve protection against endemic infectious diseases, thereby filling a need identified during the Persian Gulf War for a compact diagnostic identification tool. Today's global military mission, with ongoing operations in war-torn locations teeming with infectious diseases, demands a readily accessible, far-forward biological agent identification capability. This is critical to maintaining troop readiness, quickly determining patient treatment, disposition (for example, quarantine and medical evacuation), and protecting the homeland population from infections acquired by the military, from bioterrorism, and from emerging disease threats.\nThe DoD has addressed the biological threat in the context of the battlefield. However, biological threat reduction in the civilian population context is different. For example, the average civilian is not trained or equipped for response, the public health system is not supported with the kind of central command and control systems associated with the military, different requirements exist on sensitivity and different levels of tolerance for false positives and false negatives, and there is a need for dealing with a broader set of potential agents. Also, much higher sensitivity is required for Counter Terrorism (CT) detection, raising substantial technology challenges and the need to assess background interferences that may be more significant for low-level detection and monitoring schemes.\nThe urgent need for improved biohazard surveillance capability was also recognized and described for the first time in other Government agencies during this period. For example, the United States Postal Service has developed a Biohazard Detection System (BDS) using proven technology to implement early identification of anthrax. The BDS unit consists of an air-collection hood, a cabinet where the collection and analysis devices are housed, a local computer network connection, and a site controller (a networked computer). All BDS processes are automated. The equipment continuously collects air samples from mail canceling equipment while the canceling operation is underway. The air collection hood is installed over the canceling equipment at the very first pinch point in the mail processing operation where it absorbs and concentrates airborne particles into a sterile water base. This creates a liquid sample that is injected into a cartridge. An automated polymerase chain reaction (PCR) test is performed on the liquid sample using sophisticated DNA matching to detect the presence of anthrax (Bacillus anthracis). The test sample is compared to a template for the anthrax DNA sequence for a match. The system concentrates air samples for a one-hour period followed by the PCR test that takes approximately 30 minutes. The BDS is simultaneously concentrating particles for the next sample while the PCR test is performed for the previous sample. So while the first result requires approximately 1\u00bd hours, subsequent results are obtained every hour. Upon detection of a DNA match, the BDS computer network conveys that information to the site controller computer. Local management is notified directly by on-site BDS personnel and also by multiple forms of electronic communication from the BDS site controller. The emergency action plan is activated, the facility's building alarm is sounded and everyone in the building is evacuated. Disadvantageously, the BDS is not adapted for identifying biohazards other than the anthrax spore.\nPractitioners in the art have proposed various solutions to the sampling, detection, analysis, identification and reporting problems associated with the biohazard surveillance requirement. For example, in U.S. Pat. No. 5,895,922, Ho describes a process and apparatus for detection of viable and potentially hazardous biological particles that may be dispersed in an airstream. Ho teaches a method for directing each of the contained particles along a linear path through air, in a sequential manner, and sampling them for determination of their size, whether they are biological and viable, and whether they are present in concentrations greater than background levels. The particle size identifies the particles as respirable or not and the particles are characterized as biological and viable by subjecting each particle in turn, to 340 nm, ultraviolet laser light and looking for the emission of fluorescence, which is typically emitted from bacteria or bacterial spores. Fluorescence detected in the 400\u2013540 nm range signals the presence of nicotinamide adenine dinucleotide hydrogen, which is indicative of biological activity or viability. Ho's apparatus is compact, and power-efficient because he uses a solid state, ultraviolet laser that is actuated only when the particle is passing the laser and only if it is deemed to be a biologically viable candidate, but it is disadvantageous for use in a remote automated surveillance station.\nIn U.S. Pat. No. 6,266,428, Flanigan discloses a system and method for remote detection of hazardous vapors and aerosols by means of two differential spectral signature spectra taken in the field of view at a low spectral resolution. A first linear discriminant optimized for the low spectral resolution is applied to the first spectrum to obtain a first response, and a hazardous cloud is detected automatically in accordance with the first response. A second differential spectral signature spectrum is taken in the field of view at a higher spectral resolution and a second linear discriminant optimized for the higher spectral resolution is applied to the second spectrum to obtain a second response, which is formed into a false-color image and displayed to an operator. The operator discriminates the hazardous cloud in accordance with the image. The first and second linear discriminants can be formed by linear programming. Flanigan's system is disadvantageous for use in a remote automated surveillance station.\nIn U.S. Pat. No. 6,317,080, Baxter discloses a method of tracking airborne substances including the steps of detecting the presence of one or more airborne substances and releasing a tracking balloon into the path of the one or more airborne substances, the tracking balloon having a transmission means and a global positioning means adapted to communicate the latitude and longitude coordinates of the tracking balloon whereby the latitude and longitude coordinates of the tracking balloon are representative of the latitude and longitude of the one or more airborne substances previously detected. Baxter neither considers nor suggests solutions to the remote automated surveillance problem.\nIn U.S. Pat. No. 6,490,530, Wyatt discloses an aerosol hazard classification and early warning network that includes a large number of remote detector and analysis units, which are deployed throughout a region under surveillance for a potentially hazardous aerosol intrusion. Such aerosol threats may originate from fires, volcanic eruptions, or overt releases of biological and chemical agents dispersed in aerosol form. Among the former are the characteristic toxic aerosols released during refinery fires or explosions. The latter biological agents include bacterial spores, lyophilized bacterial cells, and virus preparations, whereas chemical agents might include various forms of nerve gasses and other anti-personnel gasses such as mustard, all commonly deployed in aerosol form. Each detector station contains an aerosol handling unit that samples and transfers ambient aerosol particles one-at-a-time through a light scattering chamber where each such particle is constrained to pass through a fine laser beam producing, thereby, an outgoing scattered light wave. The scattering chamber contains a plurality of scattered light detectors arranged to accept light scattered into different angular locations. The signals detected at each detector position are processed by a corresponding digital signal processing chip with the resulting set of digitized signals being transferred to an on-board central processing unit (CPU). The CPU analyzes the set of light scattering signals and identifies or otherwise characterizes each particle. The classification data are then stored and, on preprogrammed command, telemetered to a remote \u201ccentral station\u201d by means of an on-board telemetry unit. The central station analyzes the sets of data received from all the detector stations and then instructs, as necessary, selected detector stations via telemetric means to change their sampling and telemetry rates. As soon as sufficient data are available, the central determines the presence, threat, extent, and progress of the aerosol cloud. These factors are then telemetrically transmitted by means of alarms and warnings sent to potentially threatened regions. Although Wyatt's system is well-adapted to remote surveillance and he teaches the use of fluorescence to identify biological compounds, his \u201clight scattering\u201d data are adapted to characterizing and counting particles in an aerosol and Wyatt doesn't consider the rapid and automated identification of biological particles.\nIn U.S. Pat. No. 6,532,067, Chang, et al. describe a method for fluorescence probing of particles flowing in a fluid, including steps of defining a trigger volume in the fluid by intersecting a plurality of substantially orthogonally aimed trigger laser beams, each of a different wavelength, detecting light scattered from the vicinity of the trigger volume by a plurality of particle detectors each sensitive to a wavelength corresponding to the wavelength of a trigger laser beam, probing the particles with a pulsed laser triggered by the particle detectors, collecting fluorescence emitted from the particle in a detection volume and focusing it in a detection region, detecting the fluorescence focused in the detection region. Chang, et al. neither consider nor suggest solutions to the remote automated surveillance problem.\nIn U.S. Pat. No. 6,613,571, Cordery et al. disclose a method and system for detecting biological and chemical hazards in mail using predetermined descriptions of the hazards sought but they neither consider nor suggest solutions to the remote automated surveillance problem. In U.S. Pat. No. 6,656,253, Willey, et al. disclose a dynamic electrostatic filter apparatus for purifying air using electrically charged liquid droplets but they neither consider nor suggest solutions to the remote automated surveillance problem. In U.S. Pat. No. 6,664,550, Rader et al. describe an aerosol lab-on-a-chip (ALOC) that integrates one or more of a variety of particle collection, classification, concentration (enrichment), an characterization processes onto a single substrate or layered stack of such substrates. By mounting a UV laser diode laser light source on the substrate, or substrates tack, so that it is located down-stream of the sample inlet port and at right angle the sample particle stream, the UV light source can illuminate individual particles in the stream to induce a fluorescence response in those particles having a fluorescent signature such as biological particles, some of said particles. An illuminated particle having a fluorescent signal above a threshold signal may trigger a sorter module that separates that particle from the particle stream. But Cordery et al. consider the process control and particle stream separation problems and neither consider nor suggest solutions to the remote automated surveillance problem.\nIn view of the recent terrorism-related security requirements mentioned above, there is a clearly-felt need in the art for a robust (military-hardened) miniaturized remote system for the initial detection, localized analysis and reporting of the presence of biohazards. Such a system requires a large number of permanently-deployed remote surveillance stations each of which can operate independently and without human intervention. Such stations must be adapted for accepting updated detection information from a remote control center to permit adaptation to global changes in the threat environment, for example.\nThese unresolved problems and deficiencies are clearly felt in the art and are solved by this invention in the manner described below."}
-{"text": "1. Field of the Invention\nThis invention relates to a display or printing device which is capable, by selective energization of a number of segments, of displaying Arabic letters, numerals and words as well as European numerals.\nThere is a strong need for such a device in the manufacture of such apparatus as small computers, portable testing or communication equipment, etc. where production is suited to markets of the Arab people and where the requirements of relatively small display area and/or low production cost make the use of said device highly desirable.\n2. Description of the Prior Art\nThe main problem in developing such a device is that it is more difficult to devise an arrangement of a limited number of segments which can be selectively operated to display Arabic words than it is for languages using the Latin alphabet. In the case of the latter said languages the starburst method has been successfuly and commonly employed (see, for example, GB No. 2001468) as disclosed, for example, in United Kingdom Pat. No. GB 2,001,468 B. Among the reasons for the aforementioned difficulty is that Arabic alphabetical characters do not lend themselves to simple goemetrical representation as readily as Latin letters. Another reason for the difficulty is that most Arabic letters take on two and sometimes three character forms, depending on their position in the printed word."}
-{"text": "The present invention relates in general to techniques for obtaining improved images by compensation methods. In particular it relates to a method for compensating defective pixels in a Spatial Light Modulator (SLM), used in optical lithography. It also relates to an apparatus for patterning a work piece comprising such a method and a method for detecting defective pixels.\nLithographic production is useful for integrated circuits, masks, reticles, flat panel displays, micro-mechanical or micro-optical devices and packaging devices, e.g. lead frames and MCM\"\"s. Lithographic production may involve an optical system to image a master pattern from a computer-controlled reticle onto a workpiece. A suitable workpiece may comprise a layer sensitive to electromagnetic radiation, for example visible or non-visible light. An example of such a system is described in WO 9945435 with the same inventor and applicant as the present invention.\nSaid computer controlled reticle may be a Spatial Light Modulator (SLM) comprising a one or two dimensional array or matrix of reflective movable micro mirrors, a one or two dimensional array or matrix of transmissive LCD crystals, or other similar programmable one or two dimensional arrays based on gratings effects, interference effects or mechanical elements (e.g., shutters).\nIn general, these computer controlled reticles may be used for the formation of images in a variety of ways. These reticles, such as an SLM, include many modulating elements or pixels, in some instances million or more pixels. For example a problem with Spatial Light Modulators is that one or a plurality of pixels in a given SLM may be defective, i.e. they may not respond to a control signal as intended.\nThese defective pixels in a computer controlled reticle may limiting the resolution and accuracy available for their use in optical imaging; e.g., the production of printed patterns on a workpiece may be limited as to its line widths and accuracy.\nTherefore, there is a need in the art for a method, which effectively and precisely finds and compensates for defective pixels in the SLM.\nIn view of the foregoing background, the compensation for defective pixels in the SLM, such as for example a mirror elements stuck in a specific position, is useful to form images having sub micron line widths with tolerances approaching 5 nm.\nAccordingly, it is an object of the present invention to improve the images formed using spatial light modulators by providing an improved method for the compensation of defective pixels.\nIn a first embodiment, the invention provides a method for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation. Said method comprising the actions of providing a source for emitting electromagnetic radiation, illuminating by said radiation said SLM having a plurality of modulating elements (pixels), projecting in a writing pass an image of said modulator on said work piece, and performing a compensation for defective pixels in at least one other writing pass.\nIn another embodiment of the invention said electromagnetic radiation is a pulsed laser source.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a single compensating pixel in another writing pass.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a plurality of compensating pixel in another writing pass.\nIn another embodiment of the invention said SLM is illuminated by a radiation dose in the different writing passes.\nIn another embodiment of the invention said SLM is illuminated by different radiation intensities in the different writing passes.\nIn another embodiment of the invention said SLM is a transmissive Spatial Light Modulator.\nIn another embodiment of the invention said SLM is a reflective Spatial Light Modulator.\nIn another embodiment of the invention the pixels in said SLM is operated in an analog mainer.\nThe invention relates also to a method for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation. Said method comprising the actions of, providing a source for emitting electromagnetic radiation, illuminating by said radiation said SLM having a plurality of modulating elements (pixels), projecting an image of said SLM on the detector arrangement for measuring a dose of radiation, and performing a compensation of said defective pixel by at least one of the most adjacent pixels in the SLM.\nIn another embodiment of the invention said compensation is performed by assigning said at least one of the most adjacent pixels by a value given by subtraction of an intended pixel value by a actual pixel value divided by the number of most adjacent pixels used for compensation.\nThe invention relates also to a method for compensating the impact of at least one defective pixel in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered at least partially with a layer sensitive to electromagnetic radiation. Said method comprising the actions of, setting the pixels in said SLM in a predetermined state, illuminating by a radiation source said SLM, projecting the image of the SLM onto the detector arrangement that measures dose of the SLM pixels, identifying defective pixels, and performing a compensation for said defective pixels in at least one writing pass.\nThe invention relates also to a method for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation. Said method comprising the actions of, providing a source for emitting electromagnetic radiation, illuminating by said radiation said SLM having a plurality of modulating elements (pixels), projecting in a first writing pass an image of said modulator on said work piece using a first set of pixels in said SLM, performing a pre compensation for defective pixels in at least one subsequent writing pass in at least one prior writing pass, and projecting in at least a second writing pass said image of said modulator on said work piece using at least a second set of pixels in said SLM.\nIn another embodiment of the invention, said method further comprising the action of, performing a post compensation for defective pixels in at least one prior writing step in at least one subsequent writing pass.\nIn another embodiment of the invention a post compensation for defective pixels in at least one prior writing step in at least one subsequent writing pass is performed instead of said pre compensation.\nIn another embodiment of the invention said electromagnetic radiation is a pulsed laser source.\nIn another embodiment of the invention, said method further comprising the action of, including at least one pixel in said first set of pixels in said at least second set of pixels.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a single compensating pixel in another writing pass.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a plurality of compensating pixels in another writing pass.\nIn another embodiment of the invention said SLM is illuminated by the same radiation dose in different writing passes.\nIn another embodiment of the invention said SLM is illuminated by different radiation dose in different writing passes.\nIn another embodiment of the invention said SLM is a transmissive Spatial Light Modulator.\nIn another embodiment of the invention said SLM is a reflective Spatial Light Modulator.\nIn another embodiment of the invention the pixels in said SLM is operated in an analog manner.\nIn another embodiment of the invention an image of said pixels in said first writing pass is displaced in said SLM relative said image of said pixels in said second writing pass by one or a plurality of pixels.\nIn another embodiment of the invention an image of said pixels in said first writing pass is displaced on said workpiece relative said image of said pixels in said second writing pass by at least a fraction of a pixel.\nIn another embodiment of the invention said first set of pixels belongs to a first SLM and said second set of pixels belong to a second SLM.\nIn another embodiment of the invention said first and second SLMs are illuminated simultaneously.\nIn another embodiment of the invention said first and second SLMs are illuminated by different radiation intensities.\nThe invention relates also to an apparatus for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation, comprising a source for emitting electromagnetic radiation, a projection system to project in a first writing pass an image of said modulator on said work piece using a first set of pixels in said SLM, means for performing a pre compensation of defective pixels in at least one subsequent writing pass in at least one prior writing pass, a projection system to project in at least a second writing pass said image of said modulator on said work piece using at least a second set of pixels in said SLM, means for performing a post compensation to of defective pixels in at least one prior writing pass in at least one latter writing pass.\nIn another embodiment of the invention said electromagnetic radiation is a pulsed laser source.\nIn another embodiment of the invention at least one pixel in said first set of pixels is included in said at least a second set of pixels.\nIn another embodiment of the invention said projection system to project in at least a second writing pass comprises, said SLM reprogrammed with the image to be written on said work piece with said at least second set of pixels, a movable stage upon which stage said work piece is arranged in order to illuminate the same feature on said work piece with said at least second writing pass as said first writing pass.\nIn another embodiment of the invention said movable stage is moved the length of N SLM pixels.\nIn another embodiment of the invention said stage is moved along a row of pixels.\nIn another embodiment of the invention said movable stage is moved along a column of pixels.\nIn another embodiment of the invention said movable stage is moved along both a row of pixels and a column of pixels.\nIn another embodiment of the invention said movable stage is moved the length of N SLM pixels plus a fraction of a SLM pixel.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a single compensating pixel in another writing pass.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a plurality of compensating pixel in another writing pass.\nIn another embodiment of the invention said SLM is illuminated by a same radiation dose in the different writing passes.\nIn another embodiment of the invention said SLM is illuminated by different radiation intensities in the different writing passes.\nIn another embodiment of the invention said SLM is a transmissive Spatial Light Modulator.\nIn another embodiment of the invention said SLM is a reflective Spatial Light Modulator.\nIn another embodiment of the invention the pixels in said SLM is operated in an analog manner.\nThe invention relates also to an apparatus for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation, comprising a source for emitting electromagnetic radiation, a projection system for illuminating said SLM, having a plurality of modulating elements (pixels), by said radiation and projecting in a writing pass an image of said modulator on said work piece, the detector arrangement (65) for measuring the dose of pixels from the image of the SLM and a computer (66) for performing a compensation for defective pixels in at least one other writing pass out of said image on said detector (65).\nThe invention relates also to an apparatus for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) (30) when creating a pattern of the SLM (30) on a work piece (60) covered with a layer sensitive to electromagnetic radiation, comprising a source for emitting electromagnetic radiation, a projection system for illuminating said SLM (30), having a plurality of modulating elements (pixels), by said radiation and projecting in a writing pass an image of said modulator (30) on said work piece (60), the detector arrangement (65) for measuring the dose of pixels from the image of the SLM, and a computer (66) for performing a compensation for defective pixels (110) by using at least one of the most adjacent pixels (111, 112, 113, 114, 115, 116, 117, 118) to said defective pixel (110).\nIn another embodiment of the invention the pixel intensities are detected by said detector (65) whenever a new work piece (60) is to be patterned.\nThe invention also relates to a method for detecting at least one defective pixel in at least one SLM. Said method comprising the actions of addressing all pixels in said at least one SLM with a first steering signal, illuminating said at least one SLM with electromagnetic radiation, recording an image of said at least one SLM, computing a gradient field of the recorded image, computing a divergence of the gradient field, identifying extreme values from the computed divergence which corresponds to defective pixels.\nIn another embodiment said invention further comprising the actions of addressing all pixels said at least one SLM with a second steering signal, illuminating said at least one SLM with electromagnetic radiation, recording an image of said at least one SLM, computing a gradient field of the recorded image, computing a divergence of the gradient field, identifying extreme values from the computed divergence, where defective pixels corresponds to extreme values from said first steering signal and said second steering signal representing same pixels.\nThe invention also relates to a method for detecting at least one defective pixel in at least one SLM. Said method comprising the actions of addressing all pixels in said at least one SLM with a first steering signal, illuminating said at least one SLM with electromagnetic radiation, recording a first image of said at least one SLM, addressing all pixels in said at least one SLM with a second steering signal, illuminating said at least one SLM with electromagnetic radiation, recording a second image of said at least one SLM, computing the difference between said first image and said second image, identifying bad pixels where the computed difference has a local minimum.\nThe invention also relates to a method for detecting at least one defective pixel in at least one SLM Said method comprising the actions of addressing a pattern to said at least one SLM, illuminating said SLM with electromagnetic radiation, recording a first image of said at least one SLM, comparing said recorded image with pattern data at feature edges, identifying bad pixels where the feature edge is moved a predetermined distance.\nIn another embodiment said pattern is a chessboard pattern.\nIn another embodiment said pattern is a pattern with parallel lines.\nIn still another embodiment said method further comprising the action of addressing said pattern with another set of pixels in said at least one SLM, illuminating said SLM with electromagnetic radiation, recording a second image of said at least one SLM, comparing said recorded second image with pattern data at feature edges, identifying bad pixels where the feature edge is moved a predetermined distance.\nIn still another embodiment said method further comprising the action of comparing the feature edge movement in said first image with said second image for identifying bad pixels stuck at intermediate values.\nIn yet another embodiment said different writing passes is performed by means of one SLM.\nIn yet another embodiment said different areas of said SLM are used in the different writing passes.\nIn yet another embodiment said different writing passes is performed by means of a plurality of SLMs.\nIn yet another embodiment said different writing passes is performed by means of one SLM.\nIn yet another embodiment said different areas of said SLM are used in the different writing passes.\nIn yet another embodiment said different writing passes is performed by means of a plurality of SLMs."}
-{"text": "Computers and mobile devices, such as cellular phones and personal digital assistants, have become increasingly interconnected due to the widespread availability of wired and wireless connections to communications networks such as the Internet. Even in the earliest days of the ARPANET, users took advantage of such interconnectivity to communicate with one another through early forms of email. As email grew in availability and popularity, email \u201clists\u201d became a popular tool for broadcasting messages to predefined groups of recipients.\nIn the 1980's, Internet based \u201cnewsgroups\u201d emerged in which users could read and respond to discussion threads revolving around a wide variety of predefined categories. Newsgroups are generally readable and updatable by anyone with the equipment to access them, since access to newsgroups is generally not restricted on a per-newsgroup or per-user basis. World Wide Web based discussion groups (i.e., also known as clubs) have also provided a way for groups of people to associate around a topic. Through the use of web server programming, the idea of discussion groups and discussion threads has been extended to provide users with the ability to subscribe to secured discussion forums that are, in some cases, moderated by other users.\nAnother variant of Internet based communication forums are the web-based \u201csocial network\u201d applications, in which a number of users are able to find each others' accounts and voluntarily become \u201cfriends\u201d or \u201cfollowers\u201d of each other's posted messages. Users generally post brief messages about their status, mood, activities, and such, and their friends and followers can read and optionally reply to those messages. As such, friends may stay abreast of each other's' activities as a tool for maintaining their social bonds."}
-{"text": "Medicament delivery devices are routinely used by persons without formal medical training, i.e. patients where self-management of their condition is increasingly common. These circumstances set a number of requirements for medicament delivery devices of this kind The injector should be robust in construction, yet easy to use in terms of its operation by a user and the manipulation of the parts. In the case of those with diabetes, many users will be of impaired vision and may also be physically infirm. Devices that are too large or cumbersome may therefore prove difficult to use, particularly someone with reduced dexterity.\nU.S. Pat. No. 5,989,221 describes an electronically controlled injection device in which the readying of the device for administering and the subsequent drug delivery therefrom is controlled by an electronic control unit. Specifically, the control unit comprises a position or attitude sensor for transmitting a signal so that the readying of the device cannot take place unless the longitudinal axis of the injection cartridge is in a predetermined direction. This is in the context of removing air or mixing drug components. The control unit is also operative for driving a piston rod a predetermined distance for delivery of a drug dose.\nIn such prior art devices, the user is still required to prime the device after installation of the medicament cartridge. Moreover, such devices may under dose particularly in cases where the medicament cartridge is not properly seated within the device after insertion by the user."}
-{"text": "The invention relates to miter saws and in particular to a miter saw and a fence assembly for supporting a workpiece during cutting of the workpiece.\nCompound miter saws, typically include a base having a support surface for supporting a workpiece and a turntable mounted on the base for pivotal movement about a vertical axis. A saw arm is supported by the turntable for movement therewith, and a fence is supported by the base to provide support for a workpiece on the support surface during cutting of the workpiece. The saw arm has a saw blade and is pivotable about a horizontal axis between non-cutting and cutting positions. If it is desired to engage in compound miter cutting of the workpiece, the saw arm is pivoted about a horizontal bevel angle axis to allow the user to make the angular cuts in the workpiece. In some prior art miter saws, to provide sufficient clearance for the blade of the saw arm and yet provide adequate vertical support to the workpiece during the cutting of the workpiece, the end of the fence adjacent the saw blade has an angular edge portion. The angular edge portion is intended to provide support for the workpiece while also providing clearance for the saw blade during compound miter cutting of the workpiece. It is known in the art to adjust the position of the fence toward and away from the saw blade to allow cutting at extreme bevel angles and yet provide support to the workpiece."}
-{"text": "A mobile agent can be written as a program that executes on a set of network hosts. The agent visits the network hosts to execute parts of its program. The agent may need, for example, to access information located on a given network host or there may be some preference to execute parts of its program on various network hosts.\nIn prior art systems, the sequence of hosts that the agent visits is statically pre-configured when an agent program is written. Furthermore, the agent in these systems must execute an explicit instruction to move to another remote network host. For example, U.S. Pat. No. 5,603,031, issued Feb. 11, 1997, entitled \"System and Method for Distributed Computation Based Upon the Movement, Execution, and Interaction of Processes in a Network,\" by White et al., describes a method for statically pre-configuring an agent's itinerary in a destination list composed of destination objects. Each destination object has a telename and a teleaddress preassigned to specific regions of the network. In the system disclosed in White et al., an agent must execute a \"go\" statement to move to another network host to access resources located at that host.\nThere are many other examples of this explicit agent relocation requirement. In Lange et. al., IBM Aglets Workbench--Programming Mobile Agents in Java, Proceedings of 1997 World Wide Computing and Its Applications, Japan, pp. 253-266, the authors suggest that an agent execute the \"dispatch\" statement to move to another network host. In Cardelli, A Language with Distributed Scope, Computing Systems, Vol. 8, No. 1, Winter 1995, pp. 27-59, Cardelli describes a system wherein object migration is achieved by explicitly copying the object (agent) state. Finally, in Chess et al., Itinerant Agents for Mobile Computing, available as IBM Research Report RC-20010, the authors state the need for explicit primitives and mobility support to make an agent move to another network host.\nSystems that require this explicit agent relocation suffer from increased programming complexity because the programmer must be aware of the location where a particular piece of code would execute. As a result, the application code has to be organized into location-sensitive blocks with explicit instructions for agent relocation demarcating the blocks. Location awareness distracts the programmer from focusing solely on developing application logic. This reduces programmer productivity.\nThus, there is a need for a system which permits a program (agent) to execute throughout a network without the need for explicit agent relocation, thereby improving programmer productivity."}
-{"text": "Technical Field\nThe present invention relates to an electric power supply system.\nBackground Art\nHeretofore, there has been proposed an electric power supply system in which a direct current power supply such as a battery and fuel cell mounted on an electric vehicle such as an electric automobile and a fuel cell automobile is used to supply electricity to household electrical devices (for example, refer to Japanese Unexamined Patent Application, First Publication No. 2006-325392).\nThe electric power supply system disclosed in Japanese Unexamined Patent Application, First Publication No. 2006-325392 comprises: a vehicle having a means for supplying electric power to the outside of the vehicle; a stationary fuel cell system provided with an inverter; a load device that receives electric power supply from the stationary fuel cell system; and a system power supply that supplies electric power to the stationary fuel cell system. In the event of a power outage, this electric power supply system connects the vehicle and the stationary fuel cell system, and supplies electric power from the vehicle to the load device via the inverter of the stationary fuel cell system.\nAs a fuel cell, there is known a fuel cell stack in which a membrane electrode assembly is formed by arranging an anode electrode and a cathode electrode on either side of a solid polymer electrolyte membrane (hereunder, referred to as \u201celectrolyte membrane\u201d), arranging a pair of separators on either side of this membrane electrode assembly to form a flat unit fuel cell (hereunder, referred to as \u201cunit cell\u201d), and then stacking a plurality of these unit cells together to form a fuel cell stack. In this fuel cell, hydrogen ions produced by a catalytic reaction at the anode pass through the electrolyte membrane and move toward the cathode. There, they react with the oxygen in the air, giving rise to an electrochemical reaction and the generation of electric power.\nThe fuel cell described above generates heat as electric power generation is performed, and therefore, the generated water produced as a result of the electric power generation in the fuel cell is likely to vaporize. The generated water that has vaporized (water vapor) is discharged together with cathode off-gas and anode off-gas, and as a result, the electrolyte membrane of the fuel cell becomes dry. If the fuel cell becomes excessively dry (hereunder, referred to as \u201cdry-up condition\u201d), there is a problem in that the power generation performance of the fuel cell becomes reduced, and this consequently leads to deterioration in the electrolyte membrane.\nTherefore, the fuel cell system is provided with a cooling device for cooling the fuel cell which generates heat as power generation is performed. The cooling device is formed with a coolant that circulates in the fuel cell and absorbs heat, a radiator for releasing heat from the coolant, and a radiator fan that blows air to the radiator.\nIncidentally, in general, cooling devices of fuel cells and control thereof are designed in consideration of a vehicle in a state of traveling.\nSpecifically, they are designed such that in those cases where the traveling speed of the vehicle is high and the amount of electric power being generated by the fuel cell is high, traveling air stream is introduced into the radiator and the radiator fan is rotated at a high rotation speed, to release the heat of the coolant flowing through the radiator. Moreover, they are designed such that in those cases where the traveling speed of the vehicle is low and the amount of electric power being generated by the fuel cell is low, traveling air stream is introduced into the radiator and the radiator fan is rotated at a low rotation speed, to release the heat of the coolant flowing through the radiator. As a result of this, the fuel cell is appropriately cooled to suit the traveling state of the vehicle (to suit the electric power generation state of the fuel cell), and it is therefore possible to prevent the electrolyte membrane from becoming dry when the vehicle is traveling.\nFurthermore, the fuel cell system is designed so that rotation of the radiator fan stops when the vehicle is stopped and the electric power generation of the fuel cell is stopped. As a result, wasteful electric power consumption by the radiator fan can be prevented."}
-{"text": "Adhesives for optical films require durability so that foaming, lifting, peeling, etc. of optical films should not occur even in a high-temperature high-humidity atmosphere when optical members are bonded to adherends, and require light leakage prevention property so that the adhesives may flexibly follow dimensional changes of optical films in a high-temperature atmosphere so as not to cause light leakage in the case where two polarizing films that are optical films are bonded with the adhesive in such a manner that their polarization axes are at right angles to each other and allowed to stand at high temperature or at high temperature and high humidity. Further, rework property is also required in order that an optical member (member having an adhesive on an optical film) may be re-released from an adherend such as a liquid crystal panel without contaminating the adherend even in the case where the optical film is bonded to the adherend at the deviated position.\nAs such adhesives for optical films, acrylic-based adhesives have been mainly used in the past, and as an adhesive satisfying the above requirements, an acrylic-based adhesive obtained by blending an acrylic polymer of high-molecular weight with a medium- or low-molecular weight substance and crosslinking the resulting blend is known. In this acrylic-based adhesive, by crosslinking the acrylic polymer of high-molecular weight, cohesive force is enhanced to impart durability so as not to cause lifting and peeling, and by introducing the medium- or low-molecular weight substance, occurrence of light leakage is inhibited.\nHowever, with recent increase of sizes of displays, light leakage attributable to dimensional changes of optical films, particularly polarizing films, in a high-temperature atmosphere becomes a more serious problem, and the adhesives for optical films have been desired to have higher flexibility. In severe conditions, further, there occurs a problem that the low-molecular weight substance added for enhancing light leakage prevention property bleeds out to contaminate the adherend in the re-release process or to cause peeling.\nFor example, technique to increase a ratio (Mw/Mn) of the weight-average molecular weight (Mw) to the number-average molecular weight (Mn) to not less than 4 in order to improve heat durability and stress relaxation property has been disclosed (patent document 1). However, by merely increasing the Mw/Mn ratio to not less than 4, the stress relaxation property required for large displays is insufficient.\nIt is also known to enhance adhesion to plastics in a high-temperature high-humidity atmosphere by using a mixture of a (meth)acrylic-based polymer having a weight-average molecular weight of not less than 200,000 and a (meth)acrylic-based polymer having a weight-average molecular weight of less than 200,000 (patent document 2). In this method, however, the component having a lower weight-average molecular weight bleeds on the surface, so that problems of contamination in the re-release process and peeling under severe working conditions could not been solved.\nFurther, it is also known to user as an adhesive for optical films, a polymer in which an aromatic group-containing monomer is copolymerized (patent document 3 and patent document 4). However, its object is to control the refractive index of the adhesive in order to inhibit interfacial reflection between an optical member, such as a polarizing plate or a retardation plate, and the adhesive and between the adhesive and an adherend. Since the adhesive is designed so as to exhibit strong adhesion immediately after bonding, there is a problem that the adhesive is poor in a property that re-bonding is possible even if wrong bonding is made (rework property).\nMoreover, it is also known to use, as an adhesive for lowly polar films, a high-molecular weight polymer containing an alicyclic monomer or an aromatic group-containing monomer as a copolymer component (patent document 5). The adhesive, however, is designed so that the value (Mw/Mn) obtained by dividing the weight-average molecular weight of the copolymer by the number-average molecular weight thereof may become relatively small, and therefore, when the adhesive is used for a polarizing film having a large bond area or the like, there resides a problem that a stress due to dimensional change of the polarizing film or the like cannot be completely relaxed, and hence, sufficient light leakage prevention property is not obtained. Further, even if the alicyclic monomer is copolymerized, the light leakage prevention property is insufficient.\nPatent document 1: Japanese Patent Laid-Open Publication No. 341141/2002\nPatent document 2: Japanese Patent Laid-Open Publication No. 107507/2002\nPatent document 3: Japanese Patent Laid-Open Publication No. 173656/2002\nPatent document 4: Japanese Patent Laid-Open Publication No. 13029/2003\nPatent document 5: Japanese Patent Laid-Open Publication No. 053976/2005\nIt is an object of the present invention to provide an acrylic-based adhesive composition for optical films, which not only relaxes a stress due to dimensional change of a member having a particularly large bond area, such as a polarizing film, a retardation film or an elliptically polarizing film, and prevents contamination in the re-release process and peeling or foaming under high-humidity high-temperature conditions but also inhibits color nonuniformity due to light leakage, and to provide an optical member using the adhesive composition."}
-{"text": "1. Field of the Invention\nThe invention relates to an automatic transmission for automotive vehicles having multiple-ratios wherein ratio changes between at least one pair of ratios is a swap-shift.\n2. Background Art\nSwap-shift transmissions for automotive vehicle powertrains are disclosed in prior art U.S. Pat. Nos. 6,292,731; 5,722,519; 5,553,694; 5,758,302; 6,370,463; and 6,577,939. Each of these patents discloses a control system for a multiple-ratio automatic transmission for automotive vehicle powertrains. The transmission includes first and second gearsets arranged in series so that the torque output element of the first gearset delivers torque to the torque input element of the second gearset. The first gearset is an overdrive gearset or an auxiliary gearset. The second gearset is a multiple-ratio gearset, which is referred to as the main gearset. In the case of the prior art patents identified above, the multiple-ratio gearset is a so-called Simpson gearset, which comprises a first planetary gear unit and a second planetary gear unit with a common sun gear.\nThe torque input element of the first gearset is connected to the turbine of a hydrokinetic torque converter driven by an engine in the powertrain. The torque output element of the second gearset is connected to vehicle traction wheels through a differential-and-axle assembly in known fashion.\nDuring acceleration of the vehicle, as the engine delivers power to the traction wheels, the overall transmission ratio can progress from an initial high torque multiplication ratio to a low torque multiplication ratio with ratio steps between the high ratio and the low ratio.\nThe transmission is characterized by a so-called swap-shift during upshifts from a second overall forward driving gear ratio to a third overall forward driving gear ratio and from the second overall forward driving gear ratio to a fifth overall forward driving gear ratio. Conversely, a swap-shift occurs during downshifts from the third overall gear ratio to the second overall gear ratio and from the fifth overall gear ratio to the second overall gear ratio.\nDuring a swap-upshift from the second overall gear ratio to the third overall gear ratio, the overdrive gearset must be downshifted while the Simpson gearset is upshifted, the shifts in the gearsets being synchronized or coordinated during the overall shift progression. Conversely, a swap-downshift from the third overall gear ratio to the second overall gear ratio requires a simultaneous upshift of the overdrive gearset and a downshift of the Simpson gearset in a synchronized fashion.\nRatio changes are controlled by a pressure operated friction clutch or brake for each gearset. In the case of a 2-3 swap-upshift from the second overall gear ratio to the third overall gear ratio, a reaction brake on the overdrive gearset must be released as a reaction brake for the Simpson gearset is applied. Conversely, on a 3-2 swap-downshift, a reaction brake for the Simpson gearset must be released in synchronism with the application of a friction brake for the overdrive gearset.\nA swap-upshift from the second overall gear ratio to the fifth overall gear ratio requires release of a reaction brake for the overdrive gearset in synchronism with engagement of a friction clutch for the Simpson gearset. The friction clutch for the Simpson gearset drivably connects together two gear elements of the Simpson gearset so that the Simpson gearset gear ratio, following the upshift, is unity.\nIn the transmission of the present disclosure, the overdrive gearset is a simple planetary gear unit with two gear ratios: a low ratio, which is unity, and a high ratio, which is an overdrive. The Simpson gearset is characterized by three forward drive gear ratios (as well as a reverse ratio). The first forward drive gear ratio has high torque multiplication, and the third forward drive gear ratio is unity. The second gear ratio is approximately midway in the torque ratio progression from the highest torque multiplication ratio to the lowest torque multiplication ratio.\nThe overdrive gearset, in combination with the three Simpson gearset gear ratios, is capable of producing an overall ratio range of six gear ratios, since each of the three Simpson gearset gear ratios can be combined with each of the two overdrive gearset gear ratios.\nPrecise synchronization is required to achieve acceptable shift quality during swap-upshifts and swap-downshifts. This synchronization should be maintained throughout the life of the transmission, notwithstanding the presence of wear of the torque transfer friction elements and changes in environmental conditions, such as temperature, lubricant viscosity changes and changes in coefficients of friction for the clutches and brakes.\nEven small errors in synchronization of the gear ratio changes for the overdrive gearset and the Simpson gearset, as the two gearsets are upshifted and downshifted during a swap-shift, will have a significant adverse effect on the overall shift quality. One of these adverse effects is referred to as a ratio \u201cflare\u201d condition. Another is referred to as a ratio \u201covershoot\u201d condition. These conditions, as well as other error conditions, can cause significant, perceptible torque disturbances at the torque output shaft for the transmission.\nA ratio \u201cflare\u201d occurs when the overdrive gearset begins its gear ratio progression during a swap-shift before the Simpson gearset begins its gear ratio progression. A ratio \u201covershoot\u201d occurs if the Simpson gearset shift progression ends before slipping of the friction element for the overdrive gearset is completed."}
-{"text": "1. Field of the Invention\nThe present invention relates to a system for an intelligent attitude and orbit control of a satellite, and more particularly to improved intelligent control apparatus and method where autonomous attitude and orbit control are performed onboard a satellite in normal and contingency modes.\n2. Description of the Conventional Art\nConventionally, attitude determination and control are performed automatically onboard a satellite. There are patents to improved the accuracy of the attitude control system using inertial sensor and the star sensor. See U.S. Pat. No. 4,617,634, entitled, \"Artificial Satellite Attitude Control System,\" dated Oct. 14, 1986. On the other hand, conventional orbit control method requires tracking of the satellite from the ground, calculation of the parameters to change the orbit, and transmitting commands to the satellite. But this method requires much financial and human resources. Also there exists potential for the human error in the process. Thus, Wertz proposed the method of maintaining an assigned orbit without control or intervention from the ground. See U.S. Pat. No. 5,528,502, entitled \"Satellite Orbit Maintenance System,\" dated Jan. 18, 1996. The present invention provides a method and apparatus to control both attitude and orbit of a satellite autonomously based on intelligent decisions onboard the satellite.\nFor controlling orbit and attitude of a satellite automatically, the position of a satellite has to be determined (Navigation technique), and controlled (Guidance technique).\nNavigation and guidance technique is being actively studied and some of them are being registered as patents.\nNavigation technique is reported in the paper, \"Autonomous Navigation\" C. Jayles, F. Alby, J. Berthios, D. Pradines, Spaceflight Dynamics Part II, Edited by Jean-Pierre Carrou, CNES, 1995. And Navigation technique is disclosed as a patent, \"Method and Apparatus for predicting the position of a satellite in a satellite based navigation systems (U.S. application Ser. No. 5,430,657).\nAlso, Guidance technique is reported in \"TOPEX/POSEIDON Autonomous Maneuver Experiment (TAME) Design and Implementation\" by T. Kia, J. Mellstrom, A. Klumpp, T. Munson, and P Vaze, Advances in the Astronautical Sciences, Guidance and Control 1997, Edited by Robert Culp and Stuart Wiens, American Astronautical Socienty Publication, pp. 41-56.\nThere are several reports and registered patents, for attitude control e.g., U.S. Pat. Nos. 5,534,965, 5,458,300 and 5,412,574.\nIt is desirable to have a system where the attitude the the orbit are maintained autonomously onboard the satellite autonomously onboard the satellite during the normal operational mode. Furthermore, in the contingency situation such as when the collision danger exists or when a sensor or an actuator fails, we require the satellite to maneuver automatically to avoid the collision or to operate without the failed sensor. Thus, an intelligent decision making process for the attitude and orbit control are highly desired quality."}
-{"text": "The field of the disclosure relates generally to data communications and more specifically, to an optical fiber serial interface module that interfaces between a terminal controller and a data bus.\nAt least some known applications include an ARINC 629 data bus that uses metal twisted pair electrical bus cables, stub cables, bus terminators and current mode couplers (CMC) mounted on heavy metallic panels. However, because of such components, the data bus is bulky, heavy and expensive. Optical communications solutions, such as those that utilize optical fiber as a communications media, are desirable due to the reduced weight, which may be advantageous in an aircraft.\nOne existing solution for implementing an optical fiber data bus incorporates glass optical fibers (GOFs). This system utilizes 850 nm wavelength transmitters and receivers that are packaged individually in a pair, called a Fiber Optic Serial Interface Module (FOSIM). The FOSIM transmitter and receiver have interface electronics to the terminal controller which transmit and receive electrical signal to and from the FOSIM in Manchester bi-phase format. In the typical aircraft application, these FOSIMs are located inside the various avionics subsystems of the aircraft that utilize the data bus for communications. Often, such avionic subsystems are referred to as Line Replaceable Units (LRUs). Inside the LRU, the FOSIMs are mounted along with the terminal controller on a multilayer 6U (full size) VME circuit card.\nHowever, GOFs may be relatively fragile and break relatively easily during installation on a vehicle, such as an airplane. Further, GOFs have a relatively small diameter, which may make optical alignment difficult. Therefore, components associated with GOFs, such as connectors and optoelectronic devices, may be relatively expensive. There is a strong desire in the aircraft production community to develop an optical data bus that uses more robust optical fiber, such as a plastic optical fiber data bus, to replace the current electrical ARINC 629 data bus for future upgrades of such aircraft, although implementations would not be limited to aircraft applications."}
-{"text": "Aging of the skin is a complex phenomenon resulting from the interaction of several intrinsic and extrinsic factors. Intrinsic aging is an inevitable, genetically programmed process. Among extrinsic influences (e.g., wind, heat, cigarette smoke, chemicals, etc.), ultraviolet radiation appears to be the single most important factor associated with aging of the skin. The effect of ultraviolet radiation on elastic tissues results in elastosis, which is the accumulation of damaged elastin, resulting in reduced elasticity and resilience.\nElastin is a critical component of extracellular matrix, and is especially abundant in tissues subject to physical deformations, such as lungs, blood vessels and skin.\nThe effect of intrinsic aging on tissue elasticity of mucosal tissues (such as vaginal, oral, or rectal mucosal tissues) and of viscero-elastic tissues (that are lining body cavities such as the respiratory track, the gastrointestinal track, the urinal and bladder track, or the reproductive track) is very similar to the effect of intrinsic skin aging. Elastin fiber production in these tissues is reduced with aging, resulting in reduced responsiveness to stimuli. In the oral cavity, such changes can contribute to a decrease in the health of the gums (leading to reduced resistance to the pressure of food processing), increased gum bleeding, loose teeth, and a general decrease in the visual health parameters of the oral cavity. In the vagina, reduced elastin fiber production could result in stiffness and reduced sexual function, and uterine prolapse is associated with reduced elasticity of the female reproductive system. Reduced elasticity of the bladder can result in urine incontinence. Reduced elasticity of vessel walls can lead to vessel breakage and bruising. In the eye, degenerative changes in elastin fibers in Brunch's membrane can be responsible for deposition of drusen and macular degeneration.\nConsequently, the reduction in elasticity of these tissues results in reduced quality of life and self esteem. Thus, it is desired to have a treatment that can prevent, retard, or reverse the intrinsic and extrinsic aging effects on tissue elasticity.\nTriglycerides are a main constituent of vegetable oil and animal fats, and they play an important role in metabolism as energy sources. However, high triglyceride levels may be associated with a higher risk for atherosclerosis, heart disease, and stroke. (Forrester, J. S., Curr. Opin. Cardiology 2001, 16: 261-264). High triglyceride levels can also increase the risk of thrombosis, which can lead to myocardial infarction (Miller G. J., Atherosclerosis, 2005, 179:213-27). Hypertriglyceridemia is also a well known cause of acute pancreatitis, which can have life-threatening complications (Bae J. H. et al., Korean J Gastroenterol. 2005, 46:475-80). Current approaches for lowering triglycerides include diet and pharmacological agents, such as fibric acid derivatives, fish-oil, and CoA reductase inhibitors (Jonkers, I., et al., Am. J. Cardiovasc. Drugs 2001, 1:455-466).\nUric acid is an end product of purine metabolism. Purines are building blocks of RNA and DNA. Most uric acid produced in the body is excreted by the kidneys. An overproduction of uric acid occurs when there is excessive breakdown of cells, which contain purines, or an inability of the kidneys to excrete uric acid.\nHyperuricemia can play a role in the development of gout as well as many degenerative diseases, such as the Metabolic syndrome, which has been linked to a number of coronary heart diseases and increased mortality (Lee, M-Sh., et al., J. Clin. Nutr. 2005, 14:285:292). Hyperuricemia is also involved in the tumor lysis syndrome (TLS), which is a life-threatening constellation of metabolic derangements arising as a consequence of the release of intracellular metabolites by tumor cells as they undergo necrosis (Zeh, H J et al., J Immunother. 2005; 28:1-9). Uric acid and triglycerides were both found to be positively associated with C-reactive protein (CRP) levels (Garcia-Lorda P., et al., International Journal of Obesity (2005) 1-7).\nThus, it is desired to have a treatment that can prevent, retard, or reverse the negative cardiovascular effects induced by high blood levels of triglycerides and uric acid.\nMalvaceae is a family of flowering plants that includes the mallows, cotton plants, okra plants, hibiscus, baobab trees, and balsa trees. The family traditionally consists of about 1,500 species in 75 genera. Malva sylvestris is a species from the Malva (mallow) genera. The leaves of Malva sylvestris, otherwise known as blue mallow, are rich in mucilage. The mucilage of M. sylvestris is made up of high molecular weight acidic polysaccharides (Classen B, et al., Planta Med 64(7): 640-44 (1988)). The leaf tea is traditionally believed to be useful as an anti-inflammatory, decongestant, humectant, expectorant, and laxative. It has also been used internally for soothing sore throats, laryngitis, tonsillitis, coughs, dryness of the lungs, and digestive upsets. Mallow is also used as a poultice for healing wounds and skin inflammations. In traditional medicine, mallow leaf tea is also used against abnormal growths of the stomach and to alleviate urinary infections (Bisset N G (ed). Malvae folium\u2014Mallow leaf. In Herbal Drugs and Phyto-pharmaceuticals (1994, CRC Press, Stuttgart, pp 313-316). Studies on irritated mucus membranes have shown that the mucilage of Malva sylvestris binds to buccal membranes and other mucus membranes of the body (Schmidgall J, et al. Planta Med 66(1): 48-53(2000)).\nCotinus coggygria extract is traditionally believed to be useful as an anti-microbial treatment, used in the form of external washes. See, e.g., US Patent Applications Nos. 2002/0132021 where the extract is mentioned to be active against E. coli, Staphylococcus aureus and S. cerevisiae, as well as having anti-cancer activity. The dried leaf and twig of Cotinus coggygria is used in Chinese traditional medicine to eliminate \u201cdampness\u201d and \u201cheat\u201d, and as an antipyretic (Huang K. C., The Pharmacology of Chinese Herbs (CRS Press, 1999, pp 193-194). A yellow/orange dye can be obtained from the root and stem and can be used for fabric dying. The leaves and bark are a good source of tannins (Grieve M. A Modern Herbal. Dover Publications, Inc. NY, 1971, pp 779-781).\nThe present invention relates to the unexpected discovery that Malva sylvestris and Cotinus coggygria extracts, when ingested, are both effective for enhancing the elasticity of the skin, urogenital, blood vessel walls, and mucosal tissues, as well as reducing triglyceride and uric acid levels."}
-{"text": "This application claims benefit of Ser. No. TO2011A000483, filed 3 Jun. 2011 in Italy and which application is incorporated herein by reference. To the extent appropriate, a claim of priority is made to the above disclosed application.\nThe present invention relates to an electric propulsion system for vehicles, comprising a first and a second electric motor, a gearbox with at least three forward gears and control means arranged to control the first and the second electric motor and to control the engagement of the gears of the gearbox.\nAn electric propulsion system for vehicles of the above-identified type is known from DE-A-101 11 137."}
-{"text": "1. Field of the Invention\nThe present invention generally relates to the production techniques of integrated circuits and, more particularly, to a dynamically switched voltage screening method for quality assurance testing of integrated circuit dies on a wafer.\n2. Description of Related Art\nIntegrated circuit (IC) chips or dies fabricated on the same wafer have a wide range of performance in critical parameter characteristics due, in a large part, to process variations. As a result, out-going quality and reliability is compromised by the inconsistent performance of IC chips from the same wafer lot. Typically, stress testing is performed on the wafer during this fabrication process to eliminate the weaker chips from entering the next phase of the production cycle. For example, one of the tests regarded as one of the most severe, is that of subjecting the devices to particularly high temperatures; typically between 85xc2x0 C. and 150xc2x0 C., for accelerating the infant mortality of chips in a wafer. This test, commonly referred to as xe2x80x9cburn-inxe2x80x9d, has the objective of stimulating the failure of those devices which have developed some defects during the fabrication process and/or during handling. However, this test does not segregate parts based solely on performance characteristics. Rather, the stress is uniformly applied to all chips of a semiconductor wafer. Additionally, this test is time consuming, equipment intensive, and costly to perform.\nAlternatively, the industry has administered a voltage stress test for implementing a screen with less degrading effects. A voltage screen is basically a higher than normal voltage applied during wafer test: that effectively causes defects to manifest as failures in IC chips during subsequent verification or operational testing. The problem associated with applying a higher than normal voltage level across a semiconductor die is that some dies processed with short channel lengths have a higher tendency to fail when exposed to higher voltages. These short channel length devices would not otherwise be failures except for their vulnerability to high voltage exposure. The failures are not related to defects, rather, to the over stressing of the short channel lengths in the die. Thus, as applied, the voltage screen can be responsible for false failures, an undesirable quality assurance test result. Nevertheless, the industry standard has been to continuously apply a voltage stress at one voltage level to all die on a semiconductor wafer.\nAs discussed by Lee and Sonoda, in xe2x80x9cTEST SYSTEM FOR NARROWING THE RANGE OF PERFORMANCE CHARACTERISTICS OF MONOLITHIC INTEGRATED CIRCUITSxe2x80x9d, IBM Technical Disclosure Bulletin, Vol. 15, No. 4, September 1972, some circuits, such as those employed with FET technology, have performance characteristics that can differ by as much as 100% due to tolerances in the threshold voltage. Because of these wide differences, there are at one end of the distribution curve chips exhibiting fast response time and high power dissipation, and at the other end of the distribution curve chips having slower circuit response time and lower power dissipation. Thus, sorting the wafer at the IC chip level during quality assurance or reliability testing would be advantageous to identifying the more vulnerable short channel length devices and exposing them to less stress, thus, eliminating false failures during voltage screening.\nIn the prior art, parts have been screened during wafer testing to classify individual IC chips at various speeds. In U.S. Pat. No. 5,196,787 issued to Ovens et al. on Mar. 23, 1993, entitled, xe2x80x9cTEST CIRCUIT FOR SCREENING PARTSxe2x80x9d, a test circuit was developed on the die to measure the DC characteristics of a device, which in turn, enabled one to estimate the AC characteristics. The AC characteristic estimations were then used to screen parts into various speed classes. However, no suggestion is made to dynamically switch or adjust the stress test levels based on the operational parameters measured.\nAnother method for determining the operational speed of an IC chip is disclosed in U.S. Pat. No. 5,099,196 issued to Longwell et al. on Mar. 24, 1992, entitled, xe2x80x9cON-CHIP INTEGRATED CIRCUIT SPEED SELECTION.xe2x80x9d By forming an oscillator in an IC semiconductor chip to generate pulses representative of the speed of other components formed in the chip, the operational speed of the oscillator (typically, a ring oscillator), and therefore, that of the other components formed in the semiconductor chip, can be determined. Again, the stress test levels are not altered in response to the operation speed measurements taken.\nIC chip segregation tests also include bit-pattern recognition on each device under test. This method is particularly useful in memory device testing. In U.S. Pat. No. 4,335,457 issued to Early on June 15, 1982, entitled, xe2x80x9cMETHOD FOR SEMICONDUCTOR MEMORY TESTINGxe2x80x9d, semiconductor memory devices are tested using a special purpose computer which employs simple test patterns to determine the weakest bits of the device, and then tests only these relatively few xe2x80x9cweak bitsxe2x80x9d along with structurally and operationally adjacent bits using highly complex test patterns to determine if the device is functioning properly. Bit pattern recognition, however, is not a stress test screen. Thus, no adjustment of stress test levels, predicated on the bit pattern results, is either taught or suggested by this prior art.\nStill, other methods may be employed to distinguish the IC chips based on variations in the operational parameters. However, independent of the method chosen, some functionally operating IC chips continue to remain vulnerable to excessive stress test screening levels due to chip-to-chip process variations.\nBearing in mind the problems and deficiencies of the prior art, it is therefore an object of the present invention to provide an apparatus and method that determines the speed of IC chips on a semiconductor wafer and adjusts the stress test levels based on the speed measured for each device.\nIt is another object of the present invention to provide an apparatus and method for effectively protecting the short channel IC chip population with a lower voltage during voltage stress testing of a semiconductor wafer.\nA further object of the invention is to increase the outgoing quality and reliability of a semiconductor die using a voltage screen without falsely rejecting the short channel die during voltage stress testing of a semiconductor wafer.\nIt is yet another object of the present invention to provide an apparatus and method for segregating IC chips capable of a higher voltage withstand level without compromising the resultant yield from the wafer lot.\nAnother object of the present invention is to increase the measure of reliability of the devices on a semiconductor wafer by assigning supply current limits as a function of device speed.\nStill other objects of the invention will in part be obvious and will in part be apparent from the specification.\nThe above and other objects and advantages, which will be apparent to one of skill in the art, are achieved in the present invention which is directed to, in a first aspect, a method for testing integrated circuit semiconductor devices comprising the steps of: providing a wafer containing a plurality of integrated semiconductor devices; measuring a desired parameter of the devices; and, applying a stress test to the devices wherein test conditions of the stress test are adjusted based on the desired parameter measurements of the devices. Measuring a desired parameter first comprises verifying functionality of at least some of the integrated semiconductor devices at a set of operating conditions. The method further comprises the steps of: verifying device functionality at nominal operating conditions after the stress test; and, classifying the devices as failed if the devices do not function properly after the stress test.\nThe measuring of a desired parameter comprises measuring the operational speed of the devices prior to applying the stress test.\nAdditionally, applying a stress test comprises applying a first voltage at a value higher than the device normal operating voltage to the devices with a first measured operational speed, and a second voltage at a value lower than the first voltage to the devices with a second measured operational speed, the first operational speed being less than the second operational speed.\nThe method also comprises assigning supply current specification limits and measuring the limits of the devices such that a first supply current specification limit is assigned to devices with a faster of the operational speeds and a second supply current specification limit is assigned to devices with a slower of the operational speeds, the first supply current limit being greater than the second supply current limit.\nThe current invention is directed to, in a second aspect, a method for testing integrated circuit semiconductor devices comprising the steps of: providing a wafer containing a plurality of integrated semiconductor devices; determining functionality of the integrated semiconductor devices at nominal operating conditions; segregating the integrated semiconductor devices on the wafer by measuring a parameter of the devices; applying a stress test to the devices wherein test conditions of the stress test are adjusted based on the segregation parameter measurements of the devices; determining functionality of the devices at nominal operating conditions after the stress test; and, classifying the devices as failed if the devices do not function properly after the stress test.\nThe wafer includes devices having short and long gate channel widths, and the step of segregating the integrated semiconductor devices comprises determining which of the devices on the wafer have short gate channel widths. Applying a stress test pursuant to this method comprises applying a voltage higher than the device normal operating voltage to the devices. The step of determining which of the devices have short gate channel widths further comprises measuring the operational speed of the devices prior to applying the stress test.\nApplying a stress test comprises applying a first voltage at a value higher than the device normal operating voltage to the devices with a first measured operational speed, and a second voltage at a value lower than the first voltage to the devices with a second measured operational speed, the first operational speed being less than the second operational speed.\nMeasuring the operational speed of the devices may comprise using flush delay time measurements or a performance sort ring oscillator to determine the operational speed of the devices.\nThe wafer is fabricated to include n-type and p-type field effect transistors, and wherein the step of measuring the operational speed of the devices comprises the steps of:\ni) measuring drain-to-source current for the n-type and p-type field effect transistors;\nii) determining a drain-to-source current sum by summing the value of the n-type drain-to-source current with the absolute value of the p-type drain-to-source current;\niii) determining a device ISUM value by dividing the drain-to-source current sum by the gate channel width of the device; and,\niv) segregating the devices based on the ISUM value wherein the devices with lower ISUM values correspond to slower operational speeds and devices with higher ISUM values correspond to faster operational speeds.\nStep (iv), segregating the devices based on the ISUM value, further comprises assigning devices with ISUM values less than or equal to 700 xcexca/xcexcm as slower operational speed devices.\nThe current invention provides, in a third aspect, an apparatus for testing a wafer containing a plurality of integrated semiconductor devices including n-type and p-type field effect transistors with short and long gate channel widths, and electrical parameters comprising:\na computer containing therein a set of instructions adapted to perform the following functions:\na) measure the short and long gate channel widths of the semiconductor devices;\nb) segregate the semiconductor devices by the short and long gate channel width measurements;\nc) apply a voltage stress to the devices wherein the voltage stress is a function of the short and long gate channel width measurements; and,\nd) perform a functionality test to the devices after the voltage stress such that devices not passing the functionality test are labeled failed devices.\nIn a fourth aspect, the current invention provides a program storage device readable by machine, tangibly embodying a program of instructions executable by the machine to perform method steps for voltage screening semiconductor devices on a wafer, the method steps comprising:\na) determining functionality of the integrated semiconductor devices at nominal operating conditions;\nb) segregating the integrated semiconductor devices on the wafer by measuring a parameter of the devices;\nc) applying a stress test to the devices such that test conditions of the stress test are adjusted based on the segregation parameter measurements of the devices;\nd) determining functionality of the devices at nominal operating conditions after the stress test; and,\ne) classifying the devices as failed if the devices do not function properly after the stress test screen.\nLastly, in a fifth aspect, the current invention provides a method for testing integrated circuit semiconductor devices comprising the steps of:\na) providing a wafer containing a plurality of integrated semiconductor devices including n-type and p-type field effect transistors with short and long gate channel widths, and operational speed parameters;\nb) determining functionality of the integrated semiconductor devices at nominal operating conditions;\nc) segregating the integrated semiconductor devices on the wafer by measuring the operational speed parameters;\nd) applying a stress test to the devices wherein test conditions of the stress test are adjusted based on the operational speed measurements of the devices, such that a first voltage at a value higher than the device normal operating voltage is applied to the devices with a first measured operational speed, and a second voltage at a value lower than the first voltage is applied to the devices with a second measured operational speed, the first operational speed being less than the second operational speed;\ne) determining functionality of the devices at nominal operating conditions after the stress test; and,\nf) classifying the devices as failed if the devices do not function properly after the stress test."}
-{"text": "Customer contact centers provide an important interface for customers/partners of an organization to contact the organization. The contact can be for a request for a product or service, for trouble reporting, service request, etc. The contact mechanism in a conventional call center is via a telephone, but it could be via a number of other electronic channels, including email, online chat, etc.\nThe contact center consists of a number of human agents, each assigned to a telecommunication device, such as a phone or a computer for conducting email or Internet chat sessions, that is connected to a central switch. Using these devices, the agents generally provide sales, customer service or technical support to the customers or prospective customers of a contact center, or of a contact center's clients. Conventionally, a contact center operation includes a switch system that connects callers to agents. In an inbound contact center, these switches route inbound callers to a particular agent in a contact center or, if multiple contact centers are deployed, to a particular contact center for further routing. When a call is received at a contact center (which can be physically distributed, e.g., the agents may or may not be in a single physical location), if a call is not answered immediately, the switch will typically place the caller on hold and then route the caller to the next agent that becomes available. This is sometimes referred to as placing the caller in a call queue. In conventional methods of routing inbound callers to agents, high business value calls can be subjected to a long wait while the low business value calls are often answered more promptly, possibly causing dissatisfaction on the part of the high business value caller.\nIn many call centers, the agents answering calls are organized into a plurality of groups or teams, with each group having primary responsibility of the calls in one or more call queues. Different agent groups often have responsibility for different goals or functions of the call center, such as generating customer leads, closing sales with prospects and servicing existing customers. Routing an inbound caller to an appropriate group or team of the call center to address the needs of that caller can be a burdensome, time-consuming process.\nThere is a need for a system and method for identifying high business value inbound callers at a call center during a time period in which inbound callers are awaiting connection to an agent. Additionally, there is a need to improve traditional methods of routing callers, such as \u201cround-robin\u201d caller routing, to improve allocation of limited call center resources to high business value inbound callers. Further, there a need to improve traditional methods of routing callers to a group or team of agents appropriate to the caller's needs from a plurality of agent groups that implement different functions or goals of the call center."}
-{"text": "1. Field of the Invention\nThe field of the invention is data processing, or, more specifically, methods, apparatus, and products for profiling an application for power consumption during execution on a compute node.\n2. Description of Related Art\nThe development of the EDVAC computer system of 1948 is often cited as the beginning of the computer era. Since that time, computer systems have evolved into extremely complicated devices. Today's computers are much more sophisticated than early systems such as the EDVAC. Computer systems typically include a combination of hardware and software components, application programs, operating systems, processors, buses, memory, input/output (\u2018I/O\u2019) devices, and so on. As advances in semiconductor processing and computer architecture push the performance of the computer higher and higher, more sophisticated computer software has evolved to take advantage of the higher performance of the hardware, resulting in computer systems today that are much more powerful than just a few years ago.\nParallel computing is an area of computer technology that has experienced advances. Parallel computing is the simultaneous execution of the same task (split up and specially adapted) on multiple processors in order to obtain results faster. Parallel computing is based on the fact that the process of solving a problem usually can be divided into smaller tasks, which may be carried out simultaneously with some coordination.\nParallel computers execute applications that include both parallel algorithms and serial algorithms. A parallel algorithm can be split up to be executed a piece at a time on many different processing devices, and then put back together again at the end to get a data processing result. Some algorithms are easy to divide up into pieces. Splitting up the job of checking all of the numbers from one to a hundred thousand to see which are primes could be done, for example, by assigning a subset of the numbers to each available processor, and then putting the list of positive results back together. In this specification, the multiple processing devices that execute the algorithms of an application are referred to as \u2018compute nodes.\u2019 A parallel computer is composed of compute nodes and other processing nodes as well, including, for example, input/output (\u2018I/O\u2019) nodes, and service nodes.\nParallel algorithms are valuable because it is faster to perform some kinds of large computing tasks via a parallel algorithm than it is via a serial (non-parallel) algorithm, because of the way modern processors work. It is far more difficult to construct a computer with a single fast processor than one with many slow processors with the same throughput. There are also certain theoretical limits to the potential speed of serial processors. On the other hand, every parallel algorithm has a serial part and so parallel algorithms have a saturation point. After that point adding more processors does not yield any more throughput but only increases the overhead and cost.\nParallel algorithms are designed also to optimize one more resource\u2014the data communications requirements among the nodes of a parallel computer. There are two ways parallel processors communicate, shared memory or message passing. Shared memory processing needs additional locking for the data and imposes the overhead of additional processor and bus cycles and also serializes some portion of the algorithm.\nMessage passing processing uses high-speed data communications networks and message buffers, but this communication adds transfer overhead on the data communications networks as well as additional memory need for message buffers and latency in the data communications among nodes. Designs of parallel computers use specially designed data communications links so that the communication overhead will be small but it is the parallel algorithm that decides the volume of the traffic.\nMany data communications network architectures are used for message passing among nodes in parallel computers. Compute nodes may be organized in a network as a \u2018torus\u2019 or \u2018mesh,\u2019 for example. Also, compute nodes may be organized in a network as a tree. A torus network connects the nodes in a three-dimensional mesh with wrap around links. Every node is connected to its six neighbors through this torus network, and each node is addressed by its x,y,z coordinate in the mesh. In such a manner, a torus network lends itself to point to point operations. In a tree network, the nodes typically are organized in a binary tree arrangement: each node has a parent and two children (although some nodes may only have zero children or one child, depending on the hardware configuration). In computers that use a torus and a tree network, the two networks typically are implemented independently of one another, with separate routing circuits, separate physical links, and separate message buffers. A tree network provides high bandwidth and low latency for certain collective operations, such as, for example, an allgather, allreduce, broadcast, scatter, and so on.\nWhen processing an application, the compute nodes typically do not utilize the nodes' hardware components uniformly for each portion of the application. For example, during a portion of the application that performs a collective operation, the compute nodes typically utilize the nodes' network components that interface with the tree network but do not utilize the components that interface with the torus network. During a portion of the application that performs mathematical operations on integers, the compute nodes typically do not need to utilize the float-point units of the nodes' processors. The manner in which the nodes' hardware components are utilized to process the different portions of the application determine the overall power consumption of the nodes while executing the application. Having information on how the compute nodes consume power while executing an application may help application developers efficiently reduce the power consumption of the application, thereby conserving valuable computing resources."}
-{"text": "Organic light emitting devices (OLEDs) are comprised of several organic layers in which one of the layers is comprised of an organic material that can be made to electroluminesce by applying a voltage across the device. C. W. Tang et al., Appl. Phys. Lett. 51, 913 (1987). Certain OLEDs have been shown to have sufficient brightness, range of color and operating lifetimes for use as a practical alternative technology to LCD-based full color flat-panel displays. S. R. Forrest, P. E. Burrows and M. E. Thompson, Laser Focus World, February 1995. Since many of the thin organic films used in such devices are transparent in the visible spectral region, they allow for the realization of a completely new type of display pixel in which red (R), green (G), and blue (B) emitting OLEDs are placed in a vertically stacked geometry to provide a simple fabrication process, a small R-G-B pixel size, and a large fill factor.\nA transparent OLED (TOLED), which represents a significant step toward realizing high resolution, independently addressable stacked R-G-B pixels, was reported in U.S. Pat. No. 5,703,436, Forrest et al. This TOLED had greater than 71% transparency when turned off and emitted light from both top and bottom device surfaces with high efficiency (approaching 1% quantum efficiency) when the device was turned on. The TOLED used transparent indium tin oxide (ITO) as the hole-injecting electrode and a Mg--Ag--ITO electrode layer for electron-injection. A device was disclosed in which the ITO side of the Mg--Ag--ITO electrode layer was used as a hole-injecting contact for a second, different color-emitting OLED stacked on top of the TOLED. Each layer in the stacked OLED (SOLED) was independently addressable and emitted its own characteristic color, red or blue. This colored emission could be transmitted through the adjacently stacked transparent, independently addressable, organic layer, the transparent contacts and the glass substrate, thus allowing the device to emit any color that could be produced by varying the relative output of the red and blue color-emitting layers.\nU.S. Pat. No. 5,703,745, Forrest et al, disclosed an integrated SOLED for which both intensity and color could be independently varied and controlled with external power supplies in a color tunable display device. U.S. Pat. No. 5,703,745, thus, illustrates a principle for achieving integrated, full color pixels that provide high image resolution, which is made possible by the compact pixel size. Furthermore, relatively low cost fabrication techniques, as compared with prior art methods, may be utilized for making such devices.\nSuch devices whose structure is based upon the use of layers of organic optoelectronic materials generally rely on a common mechanism leading to optical emission. Typically, this mechanism is based upon the radiative recombination of a trapped charge. Specifically, OLEDs are comprised of at least two thin organic layers between an anode and a cathode. The material of one of these layers is specifically chosen based on the material's ability to transport holes, a \"hole transporting layer\" (HTL), and the material of the other layer is specifically selected according to its ability to transport electrons, an \"electron transporting layer\" (ETL). With such a construction, the device can be viewed as a diode with a forward bias when the potential applied to the anode is higher than the potential applied to the cathode. Under these bias conditions, the anode injects holes (positive charge carriers) into the HTL, while the cathode injects electrons into the ETL. The portion of the luminescent medium adjacent to the anode thus forms a hole injecting and transporting zone while the portion of the luminescent medium adjacent to the cathode forms an electron injecting and transporting zone. The injected holes and electrons each migrate toward the oppositely charged electrode. When an electron and hole localize on the same molecule, a Frenkel exciton is formed. These excitons are trapped in the material which has the lowest energy. Recombination of the short-lived excitons may be visualized as an electron dropping from its conduction potential to a valence band, with relaxation occurring, under certain conditions, preferentially via a photoemissive mechanism.\nThe materials that function as the ETL or HTL of an OLED may also serve as the medium in which exciton formation and electroluminescent emission occur. Such OLEDs are referred to as having a \"single heterostructure\" (SH). Alternatively, the electroluminescent material may be present in a separate emissive layer between the HTL and the ETL in what is referred to as a \"double heterostructure\" (DH).\nIn a single heterostructure OLED, either holes are injected from the HTL into the ETL where they combine with electrons to form excitons, or electrons are injected from the ETL into the HTL where they combine with holes to form excitons. Because excitons are trapped in the material having the lowest energy gap, and commonly used ETL materials generally have smaller energy gaps than commonly used HTL materials, the emissive layer of a single heterostructure device is typically the ETL. In such an OLED, the materials used for the ETL and HTL should be chosen such that holes can be injected efficiently from the HTL into the ETL. Also, the best OLEDs are believed to have good energy level alignment between the highest occupied molecular orbital (HOMO) levels of the HTL and ETL materials.\nIn a double heterostructure OLED, holes are injected from the HTL and electrons are injected from the ETL into the separate emissive layer, where the holes and electrons combine to form excitons.\nVarious compounds have been used as HTL materials or ETL materials. HTL materials mostly consist of triaryl amines in various forms which show high hole mobilities (.about.10.sup.-3 cm.sup.2 /Vs). There is somewhat more variety in the ETLs used in OLEDs. Aluminum tris(8-hydroxyquinolate) (Alq.sub.3) is the most common ETL material, and others include oxidiazol, triazol, and triazine.\nA well documented cause of OLED failure is thermally induced deformation of the organic layers (e.g. melting, crystal formation, thermal expansion, etc.). This failure mode can be seen in the studies that have been carried out with hole transporting materials, K. Naito and A. Miura, J. Phys. Chem. (1993), 97, 6240-6248; S. Tokito, H. Tanaka, A. Okada and Y. Taga. Appl. Phys. Lett. (1996), 69, (7), 878-880; Y. Shirota, T. Kobata and N. Noma, Chem. Lett. (1989), 1145-1148; T. Noda, I. Imae, N. Noma and Y. Shirota, Adv. Mater. (1997), 9, No. 3; E. Han, L. Do, M. Fujihira, H. Inada and Y. Shirota, J. Appl. Phys. (1996), 80, (6) 3297-701; T. Noda, H. Ogawa, N. Noma and Y. Shirota, Appl. Phys. Lett. (1997), 70, (6), 699-701; S. Van Slyke, C. Chen and C. Tang, Appl. Phys. Lett. (1996), 69, 15, 2160-2162; and U.S. Pat. No. 5,061,569.\nOrganic materials that are present as a glass, as opposed to a crystalline or polycrystalline form, are desirable for use in the organic layers of an OLED, since glasses are capable of providing higher transparency as well as producing superior overall charge carrier characteristics as compared with the polycrystalline materials that are typically produced when thin films of the crystalline form of the materials are prepared. However, thermally induced deformation of the organic layers may lead to catastrophic and irreversible failure of the OLED if a glassy organic layer is heated above its T.sub.g. In addition, thermally induced deformation of a glassy organic layer may occur at temperatures lower than T.sub.g, and the rate of such deformation may be dependent on the difference between the temperature at which the deformation occurs and T.sub.g. Consequently, the lifetime of an OLED may be dependent on the T.sub.g of the organic layers even if the device is not heated above T.sub.g. As a result, there is a need for organic materials having a high T.sub.g that can be used in the organic layers of an OLED.\nThe most common hole transporting material used in the HTL of OLEDs is a biphenyl bridged diamine, N,N'-diphenyl-N,N'-bis(3-methylphenyl)-1,1-biphenyl-4,4'-diamine (TPD) having the chemical structure: ##STR1##\nThis material has a good hole mobility and efficiently transfers holes to aluminum tris (8-hydroxyquinoline) in a simple single heterostructure OLED. However, TPD has a melting point of 167.degree. C. and a glass transition temperature of 65.degree. C. If a device prepared with TPD is heated above 65.degree. C., the glass transition temperature, catastrophic and irreversible failure results. In order to increase the glass transition temperature of the HTL, several groups have explored different modifications to the basic structure of TPD, Naito et al.; Tokito et al.; Shirota et al.; Noda et al. (Adv. Mater.); Han et al.; Noda et al. (Appl. Phys. Lett.); Van Slyke et al.; and U.S. Pat. No. 5,061,569. While these studies have led to materials with T.sub.g values as high as 150.degree. C., they have not led to an understanding of why certain structural modifications increase T.sub.g, while other modifications may not affect T.sub.g at all or may even lower T.sub.g. Still other modifications may produce a material not having a glass transition temperature at all or a material not having the combination of properties that is suitable for use in an HTL. For example, replacing the amine groups of TPD with carbazole groups to produce 4,4'-di(N-carbazolo)diphenyl (CBP), having the chemical structure: ##STR2## increases the melting point to 285.degree. C. However, the material shows no glass transition. Further changes in the basic structure of TPD can increase the T.sub.g value even higher, but the materials often have poorer hole transporting properties than TPD, i.e. OLEDs made with these high temperature materials give poor device properties in OLEDs compared to TPD.\nU.S. Pat. No. 5,061,569 discloses hole transporting materials comprised of at least two tertiary amine moieties and further including an aromatic moiety containing at least two fused aromatic rings attached to the tertiary amine nitrogen atoms. Out of the large number of compounds encompassed by the broadly disclosed class of compounds recited, U.S. Pat. No. 5,061,569 fails to disclose how to select those compounds which have a high glass transition temperature. For example, the naphthyl derivatives do make stable glasses. One such molecule is 4,4'-bis[N-(1-naphthyl)-N-phenyl-amino]biphenyl (.alpha.-NPD), having the chemical structure: ##STR3##\nThe present inventors' measurements show that .alpha.-NPD has a T.sub.g of 100-105.degree. C., which is substantially higher than the T.sub.g of 65.degree. C. of TPD. This material has excellent hole conduction properties, and the T.sub.g of 100-105.degree. C. is higher than the T.sub.g of TPD of about 65.degree. C. OLEDs prepared with NPD have electrical properties very similar to those prepared with TPD. However, 4,4'-bis[N-(2-naphthyl)-N-phenyl-amino]biphenyl (.beta.-NPD), having the structure: ##STR4## has been generally understood to have a T.sub.g which is substantially lower than .alpha.-derivative. Apparently because of this purportedly low and anomalous difference between T.sub.g of the .alpha.- .beta.-derivatives, there had been no known reports of using the .beta.-derivative as the hole transporting material of an OLED.\nIt would be desirable if OLED's could be fabricated from glassy charge carrier materials having improved temperature stability, while still providing luminescent characteristics comparable to prior art compounds. As used herein, the term \"charge carrier layer\" may refer to the hole transporting layer, the electron transporting layer or the separate emissive layer of an OLED having a double heterostructure. In addition, it would be useful to have a method for selecting and preparing such glassy charge carrier materials having improved temperature stability, as characterized, in particular, by glassy charge carrier materials having a high glass transition temperature.\nIn addition, there is a general inverse correlation between the T.sub.g and the hole transporting properties of a material, i.e., materials having a high T.sub.g generally have poor hole transporting properties. Using an HTL with good hole transporting properties leads to an OLED having desirable properties such as higher quantum efficiency, lower resistance across the OLED, higher power quantum efficiency, and higher luminance. There is therefore a need for a HTL having a high hole mobility and a high glass transition temperature."}
-{"text": "Synchronization is critical for telecommunication system performance. Frequency and time (time-of-day or wall-clock) synchronization is crucial for mobile wireless networks because the radios used in these networks operate in very strict bands that need separation to avoid channel interference which reduces the call quality and network capacity. Poor synchronization has also negative impact for the handover between base stations.\nMobile handsets generally derive the accurate frequency that they transmit and receive from the base stations. If the transmission frequencies are not very closely matched between adjacent cell sites, then \u201cclicks\u201d can occur when the call is being handed over (that is, switches) between base stations. In the worst case, the call would drop because the mobile handset would not be able to immediately lock onto and acquire the new signal.\nFailure to meet the timing requirements of the relevant standards would cause performance degradation for the radio access channels. In particular, this failure could compromise cell handover (especially for travelling mobile stations) and producing excess of dropped calls.\nWith increasing interest in packet networks as a common mode of communication, packet-based synchronization solutions are in high demand as alternative to PDH (plesiochronous digital hierarchy)/SDH (synchronous digital hierarchy) and GPS based solutions. Equipment vendors and telecom service providers are looking for new packet synchronization solutions with very high accuracies beyond those attainable using the traditional packet methods like Network Time Protocol (NTP) (see Mills, D., \u201cNetwork Time Protocol (Version 3) Specification, Implementation and Analysis\u201d, IETF RFC 1305, March 1992).\nIt is desirable that such new solution are also be designed for applications (e.g., base stations) that cannot bear the cost of a GPS receiver at each node, or for which GPS signals are inaccessible (for example due to location).\nThe IEEE Standard 1588 Precision Time Protocol (PTP) (IEEE Standard for a Precision Clock Synchronization Protocol for Networked Measurement and Control Systems, IEEE 1588-2008), is the latest addition in packet timing technology. Originally designed to provide precise timing for critical industrial automation applications it is now providing the highest level of accurate frequency and time to wireless backhaul networks. The backhaul portion of the network comprises the intermediate links between the core network, or backbone, of the network and the base stations. Currently standardized in 2008 as Version 2 (IEEE-1588v2, superseding IEEE 1588-2002), PTP is designed to overcome the Network Time Protocol (NTP) latency and jitter issues, providing accuracy in the nanosecond range.\nBase stations have traditionally met synchronization requirements by locking their internal oscillators to a recovered clock from the T1/E1 TDM (time division multiplexing) backhaul connection. While Ethernet has proven to be a ubiquitous and inexpensive medium for connectivity, it has not been well-suited for applications requiring precise synchronization. When the backhaul transitions from TDM to Ethernet, the base station becomes isolated from its traditional network sync feed. New base station designs are incorporating IEEE 1588 PTP slave clocks to meet the 50 ppb (parts per billion) accuracy requirement. These PTP slave clocks in the base stations rely on access to a PTP grandmaster clock deployed in a mobile switching center (MSC). Sync and Follow_Up packets flow from the grandmaster clock to the slave clocks in the base stations.\nIEEE 1588v2 PTP is fully compatible with all Ethernet and IP networks. Additionally, the protocol is designed to enable a properly designed network to deliver frequency and phase or time with precision rivalling a GPS receiver. An IEEE 1588v2 PTP Technique for Frequency Synchronization implementation can supply FDD (frequency division duplexing) and TDD (time division duplexing) radio systems and CES-based (circuit emulation services) transport systems with the synchronization signals they require as illustrated in FIG. 1. This greatly reduces the costs of clocking all wireless base station equipment using other means.\nA primary reference is a source of time and or frequency that is traceable to international standards. A recognized standard time source is a source external to PTP that provides time and/or frequency as appropriate that is traceable to the international standards laboratories maintaining clocks that form the basis for the International Atomic Time (TAI) and Universal Coordinated Time (UTC) timescales. Examples of these are Global Positioning System (GPS), NTP, and National Institute of Standards and Technology (NIST) timeservers.\nAlthough IEEE 1588v2 PTP systems add a small amount of additional traffic to the network load, they have several advantages. First, they work in the data path, and also benefit from the redundancy and resiliency mechanisms of the network, resulting in \u201calways on\u201d operation. Next, multiple transmission paths reduce redundant clock system costs. They also use a single synchronization session for all base station traffic. IEEE 1588v2 PTP supports any generic packet-based transport (such as IP, MPLS). The protocol also features configurable synchronization packet rates for network conditions to maintain accuracy.\nThe transmission of the clock information over a packet network eliminates the need for alternative mechanisms, such as GPS or prohibitively expensive oscillators placed at the receiving nodes. This provides significant cost savings in network equipment as well as in ongoing installation and maintenance. This synchronization solution transmits dedicated timing packets, which flow along the same paths with the data packets, reducing the cost of synchronization and simplifying implementation.\nTo ensure that packet technologies (Ethernet, IP, MPLS) have the necessary attributes to be truly carrier grade, operators and vendors are introducing several key technologies for the transport of timing and synchronization over packet networks. Of these IEEE 1588v2 PTP is perhaps the most important because it provides both frequency and time distribution.\nThe grandmaster is the root timing reference in a domain and transmits synchronization information to the clocks residing in its domain. In IEEE 1588v2 PTP messages are categorized into event and general messages. All IEEE 1588 PTP messages have a common header as shown in FIG. 2.\nEvent messages are timed messages in that an accurate timestamp is generated at both transmission and receipt of each message. Event messages have to be accurately timestamped since the accuracy in transmission and receipt timestamps directly affects clock distribution accuracy. A timestamp event is generated at the time of transmission and reception of any event message. General messages are not required to be timestamped. The set of event messages consists of Sync, Delay_Req (both of which have the format shown in FIG. 4), Pdelay_Req, and Pdelay_Resp. The set of general messages consists of Announce (which has the format shown in FIG. 3), Follow_Up (which has the format shown in FIG. 5), Delay_Resp (which has the format shown in FIG. 6), Pdelay_Resp_Follow_Up, Management, and Signalling.\nThe Sync, Delay_Req, Follow_Up, and Delay_Resp messages are used to generate and communicate the timing information needed to synchronize ordinary and boundary clocks using the delay request-response mechanism. A Sync message is transmitted by a master to its slaves and either contains the exact time of its transmission or is followed by a Follow_Up message containing this time. In a two-step ordinary or boundary clock, the Follow_Up message communicates the value of the departure timestamp for a particular Sync message.\nA Delay_Req message is a request for the receiving node to return the time at which the Delay_Req message was received, using a Delay_Resp message.\nThe format of the PTP message timestamp fields is shown in FIG. 7. A timestamp is the current time of an event that is recorded by a device. IEEE 1588 PTP allows for two different types of timestamping methods, either one step or two-step. One-step clocks update time information within event messages (Sync and Delay-Req) on-the-fly, while two-step clocks convey the precise timestamps of packets in general messages (Follow_Up and Delay-Resp).\nThe basic pattern of synchronization message exchange is illustrated in FIG. 8. The message exchange pattern is as follows. The master 101 sends a Sync message (M101) to the slave 102 and notes the time T1 at which it was sent. The slave 102 receives the Sync message and notes the time of reception T2. The master 101 conveys to the slave 102 the timestamp T1 by one of two ways: 1) Embedding the timestamp T1 in the Sync message. This requires some sort of hardware processing (i.e., hardware timestamping) for highest accuracy and precision. 2) Embedding the timestamp T1 in a Follow_Up message (M102). Next, the slave 102 sends a Delay_Req message (M103) to the master 101 and notes the time T3 at which it was sent. The master 101 receives the Delay_Req message (M103) and notes the time of reception T4. The master 101 conveys to the slave 102 the timestamp T4 by embedding it in a Delay_Resp message (M104).\nAt the end of this PTP message exchange, the slave possesses all four timestamps {T1, T2, T3, T4}. These timestamps may be used to compute the offset of the slave's clock with respect to the master and the mean propagation time of messages between the two clocks. The computation of offset and propagation time assumes that the master-to-slave and slave-to-master propagation times are equal.\nThe IEEE 1588 PTP based frequency recovery technique described in Section 12.1 of IEEE 1588-2008 Standard requires estimation of the mean path delay between server (master) and client (slave) which may include accounting for path asymmetry. In order to accurately synchronize to their master, slave clocks must individually determine the network transit time of the PTP messages. The network transit time is determined indirectly by measuring round-trip time from each slave to its master.\nLike all message-based time transfer protocols, PTP time accuracy is degraded by asymmetry in the paths taken by event messages. Any asymmetry in the forward and reverse path propagation times and introduces an error into the computed value of the link delay. Asymmetry can be introduced in the physical layer, e.g., via transmission media asymmetry, by bridges and routers, and in large systems by the forward and reverse paths traversed by event messages taking different routes through the network. Systems should be configured and components selected to minimize these effects guided by the required timing accuracy. In single subnet systems with distances of a few meters, asymmetry is not usually a concern for time accuracies above a few 10 s of ns. Asymmetry is not detectable by PTP; however, if known, PTP corrects for asymmetry. If two-step clocks are used, then the network has to be designed such that the general message takes the same path as the event message through a transparent clock. Failure to do this will result in a condition where the transparent clock does not calculate path delay properly. This condition is undetectable and may introduce additional jitter and wander, but it will not break the protocol.\nOther IEEE 1588 based techniques (such as those discussed in T. Neagoe and M. Hamdi, \u201cA Hardware IEEE-1588 Implementation with Processor Frequency Control,\u201d Arrow Electronics White Paper; T. Neagoe M. Hamdi and V. Cristea, \u201cFrequency Compensated, Hardware IEEE-1588 Implementation,\u201d IEEE International Symposium on Industrial Electronics, 9-13 Jul. 2006, pp. 240-245 and S. Balasubramanian, K. R. Harris, A. Moldovansky, \u201cA Frequency Compensate Clock for Precision Synchronization Using IEEE1588 and its Applications to Ethernet\u201d IEEE-1588 Workshop, September 2003 assume that the PTP GrandMaster Clock (GMC) sends Sync messages at fixed intervals, an assumption which might not necessary hold true in practice. The above receiver clock recovery mechanisms are designed based on this assumption.\nAn object of the present invention is to provide a method and system for frequency synchronization that allows one or more receivers (slaves) to frequency synchronize to a transmitter (master). Other applications of frequency synchronization are in process and manufacturing industries like paper mills, printing presses, automation and robotic systems, test and measurement instruments and systems, etc.\nAnother object of the present invention is to provide a synchronization technique based on, for example, the IEEE 1588 Precision Time Protocol (PTP) that allows frequency to be distributed over a packet network from a PTP server (master) to synchronization clients (slaves)."}
-{"text": "The present invention relates to a fungicidal mixture which comprises\na) a phenyl benzyl ether derivative of the formula I.a, I.b or I.c, \nxe2x80x83and/or a carbamate of the formula Id \nwhere X is CH or N, n is 0, 1 or 2 and R is halogen, C1-C4-alkyl or C1-C4-haloalkyl, it being possible for the radicals R to be different if n is 2, or a salt or adduct thereof, and\nb) a N-acetonylbenzamide of the formula II \nxe2x80x83where:\nR1 and R3 independently of one another are each halogen or C1-C4-alkyl;\nR2 is cyano, C1-C4-alkyl, C2-C4-alkenyl, C2-C4-alkynyl or C1-C4-alkoxy;\nR4 is hydrogen or C1-C4-alkyl;\nR5 is C2-C4-alkyl;\nR6 is thiocyano, isothiocyano or halogen, or a salt or adduct thereof,\nin a synergistically effective amount.\nThe invention further relates to methods for controlling harmful fungi using mixtures of the compounds I (I.a, I.b and I.c) and II and to the use of the compound I and the compound II for preparing such mixtures.\nThe compounds of the formula Ia-c, their preparation and their activity against harmful fungi are known from the literature (EP-A 253 213; EP-A 254 426; EP-A 398 692).\nThe compounds of the formula Id, their preparation and their activity against harmful fungi are known from the literature (WO-A 93/15,046; WO-A 96/01,256 and WO-A 96/01,258).\nAlso known are synergistic mixtures of the compounds II with dithiocarbamates, their preparation and their activity against harmful fungi (EP-A 753 258; U.S. Pat. No. 5,304,572).\nIt is an object of the present invention to provide mixtures which have an improved activity against harmful fungi combined with a reduced total amount of active ingredients applied (synergistic mixtures), with a view to reducing the application rates and to improving the activity spectrum of the known compounds I and II.\nWe have found that this object is achieved by the mixture defined at the outset. In addition, we have found that better control of harmful fungi is possible by applying the compound I and the compound II simultaneously, separately as well as together, or by applying the compound I and the compounds II in succession, than when the individual compounds are used.\nThe formula Id in particular represents carbamates in which the combination of the substituents corresponds to a line of the Table below:\nParticular preference is given to the compounds I.12, I.23, I.32 and I.38.\nThe formula II in particular represents N-acetonylbenzamides in which the combination of the substituents corresponds to a line of the Table below:\nParticular preference is given to those N-acetonylbenzamides which are generally and particularly preferred in EP-A 753 258.\nOwing to the basic character of their nitrogen atoms, the compounds Id and II are capable of forming adducts or salts with inorganic or organic acids or with metal ions.\nExamples of inorganic acids are hydrohalic acids such as hydrofluoric acid, hydrochloric acid, hydrobromic acid and hydriodic acid, sulfuric acid, phosphoric acid and nitric acid.\nSuitable organic acids are, for example, formic acid, carbonic acid and alkanoic acids, such as acetic acid, trifluoroacetic acid, trichloroacetic acid and propionic acid, and also glycolic acid, thiocyanic acid, lactic acid, succinic acid, citric acid, benzoic acid, cinnamic acid, oxalic acid, alkylsulfonic acids (sulfonic acids having straight-chain or branched alkyl radicals of 1 to 20 carbon atoms), arylsulfonic acids or aryldisulfonic acids (aromatic radicals, such as phenyl and naphthyl, which carry one or two sulfo groups), alkylphosphonic acids (phosphonic acids having straight-chain or branched alkyl radicals of 1 to 20 carbon atoms), arylphosphonic acids or aryldiphosphonic acids (aromatic radicals, such as phenyl and naphthyl, which carry one or two phosphoric acid radicals), it being possible for the alkyl or aryl radicals to carry further substituents, eg. p-toluenesulfonic acid, dodecylbenzenesulfonic acid, salicylic acid, p-aminosalicylic acid, 2-phenoxybenzoic acid and 2-acetoxybenzoic acid, etc.\nSuitable metal ions are, in particular, the ions of the elements of the first to eighth sub-group, in particular chromium, manganese, iron, cobalt, nickel, copper, zinc and of the second main group, in particular calcium and magnesium, and of the third and fourth main group, in particular aluminum, tin and lead. The metals can exist in the various valencies which they can assume.\nWhen preparing the mixtures, it is preferred to employ the pure active ingredients I and II, to which further active ingredients against harmful fungi or other pests, such as insects, arachnids or nematodes, or else herbicidal or glineth-regulating active ingredients or fertilizers can be admixed.\nThe mixtures of the compounds I and II, or the simultaneous joint or separate use of the compounds I and II, have outstanding action against a wide range of phytopathogenic fungi, in particular from the classes of the Ascomycetes, Basidiomycetes, Phycomycetes and Deuteromycetes. Some of them act systemically and can therefore also be employed as foliar and soil-acting fungicides.\nThey are specially important for controlling a large number of fungi in a variety of crop plants, such as cotton, vegetable species (eg. cucumbers, beans, tomatoes, potatoes and cucurbits), barley, grass, oats, bananas, coffee, maize, fruit species, rice, rye, soya, grapevine, wheat, ornamentals, sugar cane, and a variety of seeds.\nThey are particularly suitable for controlling the following phytopathogenic fungi: Erysiphe graminis (powdery mildew) in cereals, Erysiphe cichoracearum and Sphaerotheca fuliginea in cucurbits, Podosphaera leucotricha in apples, Uncinula necator in grapevines, Puccinia species in cereals, Rhizoctonia species in cotton, rice and lawns, Ustilago species in cereals and sugar cane, Venturia inaequalis (scab) in apples, Helminthosporium species in cereals, Septoria nodorum in wheat, Botrytis cinerea (gray mold) in strawberries, vegetables, ornamentals and grapevines, Cercospora arachidicola in groundnuts, Pseudocercosporella herpotrichoides in wheat and barley, Pyricularia oryzae in rice, Phytophthora infestans in potatoes and tomatoes, Plasmopara viticola in grapevines, Pseudoperonospora species in hops and cucumbers, Alternaria species in vegetables and fruit, Mycosphaerella species in bananas and Fusarium and Verticillium species.\nFurthermore, they can be used in the protection of materials (eg. in the protection of wood), for example against Paecilomyces variotii. \nThe compounds I and II can be applied simultaneously, that is either together or separately, or in succession, the sequence, in the case of separate application, generally not having any effect on the result of the control measures.\nThe compounds I and II are usually applied in a weight ratio of 10:1 to 0.01:1, preferably 5:1 to 0.05:1, in particular 1:1 to 0.05:1.\nDepending on the nature of the desired effect, the application rates of the mixtures according to the invention are, in particular in agricultural crops, from 0.01 to 8 kg/ha, preferably 0.1 to 5 kg/ha, in particular 0.5 to 3.0 kg/ha.\nThe application rates of the compounds I are from 0.01 to 2.5 kg/ha, preferably 0.05 to 2.5 kg/ha, in particular 0.05 to 1.0 kg/ha.\nCorrespondingly, in the case of the compounds II, the application rates are from 0.01 to 10 kg/ha, preferably 0.05 to 5 kg/ha, in particular 0.1 to 2.0 kg/ha.\nFor seed treatment, the application rates of the mixture are generally from 0.001 to 250 g/kg of seed, preferably 0.01 to 100 g/kg, in particular 0.01 to 50 g/kg.\nIf phytopathogenic harmful fungi are to be controlled, the separate or joint application of the compounds I and II or of the mixtures of the compounds I and II is effected by spraying or dusting the seeds, the plants or the soils before or after sowing of the plants, or before or after plant emergence.\nThe fungicidal synergistic mixtures according to the invention, or the compounds I and II, can be formulated for example in the form of ready-to-spray solutions, powders and suspensions or in the form of highly concentrated aqueous, oily or other suspensions, dispersions, emulsions, oil dispersions, pastes, dusts, materials for broadcasting or granules, and applied by spraying, atomizing, dusting, broadcasting or watering. The use form depends on the intended purpose; in any case, it should guarantee as fine and uniform as possible a distribution of the mixture according to the invention.\nThe formulations are prepared in a manner known per se, eg. by adding solvents and/or carriers. It is usual to admix inert additives, such as emulsifiers or dispersants, to the formulations.\nSuitable surfactants are the alkali metal salts, alkaline earth metal salts and ammonium salts of aromatic sulfonic acids, eg. ligno-, phenol-, naphthalene- and dibutylnaphthalenesulfonic acid, and of fatty acids, alkyl- and alkylarylsulfonates, alkyl, lauryl ether and fatty alcohol sulfates, and salts of sulfated hexa-, hepta- and octadecanols, or of fatty alcohol glycol ethers, condensates of sulfonated naphthalene and its derivatives with formaldehyde, condensates of naphthalene or of the naphthalenesulfonic acids with phenol and formaldehyde, polyoxyethylene octylphenol ether, ethoxylated isooctyl-, octyl- or nonylphenol, alkylphenol polyglycol ethers, tributylphenyl polyglycol ethers, alkylaryl polyether alcohols, isotridecyl alcohol, fatty alcohol/ethylene oxide condensates, ethoxylated castor oil, polyoxyethylene alkyl ethers or polyoxypropylene [lacuna], lauryl alcohol polyglycol ether acetate, sorbitol esters, lignosulfite waste liquors or methylcellulose.\nPowders, materials for broadcasting and dusts can be prepared by mixing or jointly grinding the compounds I or II or the mixture of the compounds I and II with a solid carrier.\nGranules (eg. coated granules, impregnated granules or homogeneous granules) are usually prepared by binding the active ingredient, or active ingredients, to a solid carrier.\nFillers or solid carriers are, for example, mineral earths, such as silica gel, silicas, silica gels, silicates, talc, kaolin, limestone, lime, chalk, bole, loess, clay, dolomite, diatomaceous earth, calcium sulfate, magnesium sulfate, magnesium oxide, ground synthetic materials, and fertilizers, such as ammonium sulfate, ammonium phosphate, ammonium nitrate, ureas, and products of vegetable origin, such as cereal meal, tree bark meal, wood meal and nutshell meal, cellulose powders or other solid carriers.\nThe formulations generally comprise 0.1 to 95% by weight, preferably 0.5 to 90% by weight, of one of the compounds I or II or of the mixture of the compounds I and II. The active ingredients are employed in a purity of from 90% to 100%, preferably 95% to 100% (according to NMR spectrum or HPLC).\nThe compounds I or II, or the mixtures, or the corresponding formulations, are applied by treating the harmful fungi or the plants, seeds, soils, areas, materials or spaces to be kept free from them with a fungicidally effective amount of the mixture, or of the compounds I and II in the case of separate application.\nApplication can be effected before or after infection by the harmful fungi."}
-{"text": "1. Field of the Invention\nThe present invention relates to flexible magnetic sheets, which are applied to the skin of persons or animals for therapeutic or analgesic purposes.\n2. Background Of The Invention\nMagnetic therapeutic plasters are described in U.S. Pat. No. 4,489,711 to Latzke (\"Latzke\"), which is incorporated by reference herein. The magnetic poles in FIG. 1 of Latzke are disposed in alternating rows of \"+\" and \"-\" (i.e., North and South) poles. Latzke implemented the invention on magnetic rubber sheets having a thickness of 0.5 to 1.5 mm thick, with the magnetic poles 5-10 mm apart, and a magnetic field strength of 50 to 10,000 gauss, preferably, 400 to 2,000 gauss. The examples in Latzke describe the beneficial effects of treatment with the magnetic plasters.\nU.S. Pat. No. 4,549,532 to Baermann (\"Baermann\") discloses several different geometries for a flexible magnetic sheet, including circular strips of alternating polarities (FIG. 1), triangular strips (FIG. 2), rectangular or square strips (FIG. 3) and octagonal strips (FIG. 4). U.S. Pat. Nos. 5,277,692, 5,538,495 and 5,514,072 to Ardizzone disclose a flexible magnetic pad in which alternating N/S magnets spiral out from a central core such that triangular wedges of N/S alternating magnets radiate from the central core. U.S. Pat. No. 5,304,111 to Mitsuno et al. disclose a flexible magnetic sheet in which the N/S magnetic poles are disposed in a checkerboard pattern. Other U.S. patents of interest include U.S. Pat. Nos. 4,162,672 to Yazaki, U.S. Pat. No. 5,336,498 to Snider, U.S. Pat. No. 5,017,185 to Baermann and U.S. Pat. No. 5,259,892 to Kubota. Foreign patents of interest include French Patent No. 1,215,110 to Tanaka, Japanese Laid-Open Patent Application No. 56-7405 to Miyake, and French Patent 2,371,916 to Van Den Bulke."}
-{"text": "1. Field of the Invention\nThe present invention relates to a developer for developing latent electrostatic images for use in electrophotography, electrostatic printing and the like.\n2. Discussion of Background\nTwo-component dry developers comprising carrier particles and toner particles have been conventionally known. In such a two-component dry developer, finely-divided toner particles are held on the surface of compartively large carrier particles by the electric force generated by the friction between the carrier particles and the toner particles. When the two-component dry developer is caused to come close to a latent electrostatic image, the attraction force generated between the toner particles and the latent electrostatic images overcomes the bonding force between the toner particles and the carrier particles, so that the toner particles are caused to be deposited on the latent electrostatic images. As a result, the latent electrostatic image is developed with the toner particles to a visible toner image. Therefore, the two-component dry developer is used as the toner particles are supplied thereto from time to time while in use in compensating for the toner particles used.\nAs the materials for a carrier for two-component dry developers, metallic oxides such as magnetite and ferrite are widely used. This is because such metallic oxides have a smaller apparent density than that of an iron powder carrier, so that when such metallic oxides are used as the materials for a carrier, the weight of the two-component dry developer can be reduced. Furthermore, such metallic oxides have an advantage over other materials that when such a metallic oxide is used as a carrier for a two-component dry developer, the stirring resistance of the two-component developer in a development unit is smaller than the stirring resistance of other materials employed in the carrier.\nIn addition, such metallic oxides have a smaller residual magnetic flux density and a smaller anti-magnetization force than those of an iron powder carrier, and accordingly have a smaller hysteresis loop area than that of an iron powder carrier. Furthermore, such metallic oxides have the characteristics that initial characteristics are always maintained against magnetic reversion and magnetization hysteresis.\nSince magnetite and ferrite are oxides, they are chemically stable and hardly chemically changed in contact with ozone, NO.sub.x and the like, which are formed within a copying machine.\nA carrier comprising an oxide such as ferrite or magnetite, however, has the shortcoming that a so-called spent phenomenon that a toner film is formed on the surface of carrier particles takes place by the heat generated while in use by the collision among developer particles during high speed development or during the process of making a number of copies, or by a mechanical collision between developer particles and members for a development unit while in use. Once such a spent phenomenon takes place, the charging performance of the carrier is decreased with time while in use. As a result, the toner particles are scattered and toner particles are deposited on the background of images.\nIn order to prevent the occurrence of such a spent phenomenon, methods of coating the surface of the core particles of carrier particles with a variety of resins have been proposed. However, none of them is satisfactory for use in practice. To be more specific, carrier particles coated with styrene-methacrylate resin or styrene polymer have excellent charging characteristics, but the critical surface tension thereof is relatively high, so that the above-mentioned spent phenomenon takes place during a repeated copy making operation and the life of the developer is shortened.\nExamples of a conventional negative charge controlling agent include metal complex salts of monoazo dyes, nitrohumic acid and salts thereof, sulfonated copper phthalocyanine pigments, nitro-group- or chlorine-introduced styrene oligomers, chlorinate paraffin, and melamine resin. These compounds have a complicated structure, and therefore are unstable in the properties.\nWhen such negative charge controlling agents are kneaded with application of heat thereto, they are easily decomposed and are caused to deteriorate, so that the charge controlling performance is lowered. Furthermore, the chargeability of many of such charge controlling agents is changed by the ambient conditions thereof.\nThere is a case where when a toner comprising such a conventional charge controlling agent is used for an extended period of time, the toner is deposited in the form of a film on a photoconductor because of its improper chargeability.\nFor instance, Japanese Laid-Open Patent Application 61-223753 discloses toners comprising aromatic hydroxy metallic salts such as a salicylic acid chromium complex. These toners, however, have the shortcomings that the chargeability is unstable, and the charging performance is greatly changed depending upon the ambient conditions thereof.\nJapanese Laid-Open Patent Application 3-1162 discloses a method of using a fluorinated ammonium compound or iminium compound. However, when this method is employed, the charging stability differs depending upon a carrier employed, and it is difficult to obtain a sufficient charge-stabilizing effect on a non-coated carrier for use in practice by this method.\nIn the case of a styrene-acrylic copolymer coated carrier, the chargeability is stable in a continuous mixing process, but when it is repeatedly used with a toner being replenished in a development unit, the chargeability is unstable, and the initial charge-rising performance is not satisfactory for use in practice."}
-{"text": "One example of a conventional indexing method for storage, display and search of video data is described in Patent Document 1. With this conventional method, a determination is made whether telop characters are displayed for each frame of input video data, a character area is extracted from the frame whose telop characters are displayed so that a process for recognizing the telop characters is executed. Further, an index file which includes the recognized result of the telop characters and ID information at the time of displaying the telop characters is generated. The ID information is ID information at the time of inputting the frame to be recognized. Further, the Patent Document 1 discloses that the ID information is occasionally the input time information of the frame.\nPatent Document 1: Japanese Patent Application Laid-Open No. 11-167583 (JP-A 11-167583) (paragraph 0002 and FIG. 2)"}
-{"text": "Existing shower head are mostly switched by linkage of a decorative cover and a top cover, the top cover is disposed with a plurality of outlet holes, the decorative cover has a plurality of outlet holes as well. Chinese patent database with publishing number CN201195132Y is provided with a multi-functional back-switching type hand shower head, which is disclosed with a shower head comprising a handle, an outlet housing and a switch device, the switch device comprises a switch rotation disk and a water diversion disk, the shower head is switched by rotating the switch device to control the outlet waterways. Although it is switched from a back side, the outlet housing is still disturbed with outlet holes, thus making the switch mechanism complicated and not attractive."}
-{"text": "Metallic materials typified by iron are commonly used in a state of being coated for the purpose of protection from thermal and chemical damage factors, thereby improving durability and imparting an attractive appearance. These coated metals likely require various properties such as corrosion resistance, contamination resistance, and heat resistance. Applications utilizing heat resistance include, for example, exhaust components for vehicles and motorcycles, cooking utensils, and air conditioners and heating units. These applications generally require resistance to temperatures of about 300 to 400\u00b0 C. and, in particular, exhaust components for vehicles require resistance to temperatures of 500\u00b0 C. or more.\nFor these applications, since a postcoating method of coating after forming into a predetermined shape leads to high cost and is not suitable for mass production, a precoating method of working after previously coating has been reviewed. According to the precoating method, however, working is performed after forming a film, and thus extreme workability is generally required for the coating film. For example, a thin steel sheet likely requires the workability sufficient to cause neither cracking nor peeling of the film in a 2T to 4T bending test, such that a 180\u00b0 bending is performed in the state where two to four steel sheets each having the same thickness as that of a test material are laid one upon another, and the adhesion between the film capable of enduring the above test and an underlaying metal.\nA silicone resin is known as a material of a film for heat-resistant precoated steel sheet. The heat resistance of the coating film made of the silicone resin remarkably varies depending on the kind and content of the organic group to be introduced into the silicone resin. In general, as the content of the organic group increases, the resultant film is more flexible and is excellent in workability and adhesion on working. On the other hand, the content of the organic substance should be decreased so as to enhance the heat resistance of the film. Thus, the above-mentioned workability and adhesion on working of the film are deteriorated.\nJapanese Patent Application, First Publication Nos. S63-172640, H02-265742 and H08-10701 describe, as coated steel sheets which satisfy both heat resistance and workability using a silicone resin, precoated steel sheets produced by using a coating composition containing, as a main component, a silicone resin having an alkyl group, an alkenyl group, and a phenyl group introduced therein. These coating films made of the silicone resin have a workability that is typically required of the precoated steel sheet because of comparatively high content of the organic group, but have a low resistance to temperatures of 400\u00b0 C. or less. Japanese Patent Application, First Publication No. H08-245922 describes a coating film made of a monomethylsilicone resin having a comparatively low content of the organic group, and the resin coating film is generally insufficient in workability.\nTo solve such a problem, Japanese Patent Application, First Publication No. 2002-234109 describes a method of using two kinds of resins having different heat resistances in combination. Japanese Patent Application, First Publication No. 2002-307606 describes a resin coating film containing a methylsilicone resin and a linear methylphenylsilicone resin. To improve the workability of the above-mentioned coating film described in Japanese Patent Application, First Publication No. H08-245922, Japanese Patent Application, First Publication No. 2002-80974 describes a method of using a coating composition containing a monomethylsilicone resin as a main component, an isocyanate, and an epoxy resin. As described above, however, satisfactory heat resistance and workability imparted by introduction of an organic component conflict with each other, and there have never been obtained those which generally satisfy both properties at a high level.\nThe heat treatment temperature should preferably be set to a high temperature so as to form a film having high heat resistance, and the heat-resisting temperature of the film tends to decrease so as to decrease the heat treatment temperature. It was difficult to obtain a precoated steel sheet having high heat resistance because of this problem."}
-{"text": "The field of the invention relates to memory devices and more particularly to non-volatile semiconductor memories.\nContinuing to increase rapidly is the use of computer memory, in particular non-volatile semiconductor memory, which retains its stored information even when power is removed. A wide variety of non-volatile memories exist. A typical commercial form of non-volatile memory utilizes one or more arrays of transistor cells, each cell capable of non-volatile storage of one or more bits of data.\nNon-volatile memory is unlike volatile random access memory (\u201cRAM\u201d), which is also solid-state memory, but does not retain its stored data after power is removed. The ability to retain data without a constant source of power makes non-volatile memory well adapted for consumer devices. Such memories are well adapted to small, portable devices because they are typically relatively small, have low power consumption, operate quickly, and are relatively immune to the operating environment.\nIn general, the small size, low power consumption, high speed and immunity to environment are derived from the structure of the memory. In this regard, such non-volatile memory devices are typically fabricated on silicon substrates. In addition, to obtain the advantages of small size, etc., and well as reduce costs, there is a continual effort to fabricate more circuitry within a given area.\nHighly effective approaches to increase density of nonvolatile memory include monolithic three dimensional memories disclosed in Johnson et al. U.S. Pat. No. 6,034,882, Johnson et al. U.S. Pat. No. 6,525,953, Knall et al. U.S. Pat. No. 6,420,215, and Vyvoda et al. U.S. Pat. No. 6,952,043, all hereby incorporated by reference in the entirety for all purposes.\nThe fabrication of these high-density, three dimensional memory arrays presents a number of challenges. For instance, misalignment of features during fabrication results in reduced yield and becomes more problematic as feature size is reduced. For example, in the event that a photomask is improperly placed, a memory element may be short circuited during subsequent fabrication operations. Thus, alternate methods of fabrication are needed that reduce the difficulties of aligning memory elements during fabrication while permitting improved density, decreased future size, and improved yield."}
-{"text": "Technical Field\nThe disclosure relates to a phosphor, and in particular it relates to an illumination device utilizing the same.\nDescription of the Related Art\nThere are many manufacturing methods for making white light-emitting diodes (WLEDs), such as (1) applying a yellow phosphor onto a blue light-emitting diode (LED) chip, (2) applying a red phosphor and a green phosphor powder to a blue LED chip, (3) mixing a red LED chip, a green LED chip, and a blue LED chip, and (4) applying a blue phosphor, a green phosphor, and a red phosphor (or different phosphor powders of different colors) onto an ultra-violet (UV) LED chip.\nThe WLED has several advantages over traditional incandescent light bulbs, such as long lifespan, low power consumption, small volume, fast response time, and good shake-resistance. As a result, it has gradually replaced traditional lighting products. However, current WLEDs still need to overcome problems in development, such as heat dissipation, insufficient brightness, and their relatively high price. In the lighting market, WLEDs are mainly used in auxiliary lighting such as flash lights, car interior lights, or decorative architectural lighting products. WLEDs are expected to replace traditional lighting products in the future to enter the mainstream of the global lighting market.\nIn addition to packaging techniques, the phosphor that is chosen is an important factor in the luminous efficiency of a light source. One of the directions of the research being conducted by photo electronic semiconductor companies and lighting companies is modifying phosphor compositions to increase phosphor conversion efficiency. The color render index of the white light generated by a yellow phosphor powder excited by a traditional single blue chip is not good, as the color saturation of an object illuminated by the white light is poor, thereby lowering its value in the commercial lighting market. Compared with this method, the application of a green phosphor and a red phosphor onto a blue LED may produce a white light with a better color render index. Accordingly, different phosphor compositions for use in a white light illumination device are called for."}
-{"text": "Demand for bandwidth is driving the expansion of optical transmission systems into homes and businesses of all sizes. Single wavelength fiber optic systems can support substantial data rates. However, services such as HDTV, on-demand TV programming, internet telephony, and telepresence are bandwidth intensive beyond the capabilities of many traditional networks. The present invention relates to integrated wavelength selectable photodiodes and their application to Fiber-To-The-X (FTTX) services. Fiber-To-The-X services refer to the extension of optical data transport into areas traditionally served by electrical communications systems, such as homes and small and medium sized businesses. Examples of FTTX systems are Fiber-To-The Home (FTTH), Fiber-To-The Curb (FTTC) and Fiber-To-The-Premises. FTTX architectures are also used for some highly secure optical communications links, such as radar tower interfaces."}
-{"text": "This invention is related to hybrid circuit chips or silicon substrate circuit flip chips which are soldered to a support structure. Specifically, there is provided a heating means integral with the circuit chip to permit unsoldering a specific chip which is soldered to a support structure to permit soldering a chip to a support structure without adversely affecting a neighboring chip and/or printed electrical conductors on the support structure.\nHybrid circuits usually include a plurality of circuit chips having one or more electrical components. These circuit chips are usually provided with a plurality of bonding pads for soldering to corresponding bonding pads on a glass or ceramic substrate which provide a support structure for the circuit chips. The hybrid circuit substrate or support structure also has a plurality of electrical conductors printed thereon using, typically, thin film or thick film techniques. The electrical conductors provide electrical interconnections between the circuit chips which have been soldered to the support structure.\nMany circuit chips are silicon monolithic integrated circuits. This variety of circuit chips has bonding pads on one surface thereof and includes \"solder bumps\" on the pads. This type of chip is sometimes referred to as a flip chip. The solder bumps are part of the soldering process for soldering the circuit chip to the hybrid circuit support structure as is well known.\nSince there is a very high cost in a completely assembled hybrid circuit, it is desirable to be able to repair an assembled hybrid circuit by removing a defective circuit chip and replacing it with an operative unit rather than scrapping the whole hybrid circuit. Prior art techniques for removing circuit chips from hybrid circuits include using a hand soldering iron, and using a rework fixture employing a hydrogen flame as a heating mechanism. These techniques, in many situations, result in irreparable damage to the electrical conductor patterns on the hybrid circuits as well as to adjacent good circuit chips and components. Another technique for removing hybrid circuits is shown in U.S. Pat. No. 3,904,100 in which a fixture is used in repairing hybrid circuits. The fixture includes a mechanism for heating a specific circuit chip to melt the soldering material between the circuit chip and the hybrid circuit substrate and permit removal of the circuit chip without damaging either adjacent circuit components or the conductor pattern on the hybrid circuit substrate. The latter technique for removing hybrid circuits has a disadvantage insofar as being a cumbersome apparatus to operate and requires an undue expense for providing the intended function."}
-{"text": "This invention relates to a process for the production of alkali metal silicate organic plastics by emulsifying a polymerable unsaturated organic compound and an organic epoxide compound with an aqueous alkali metal solution by mixing the polymerable organic compound and an organic epoxide compound with an aqueous solution of alkali metal silicate then adding a salt-forming compound in the amount up to 10%, based on the alkali metal silicate, preferably an organic acid, while agitating thereby producing a stable emulsion. A polymerizing catalyst such as a peroxide type catalyst is added to the emulsion thereby producing a poly(alkali metal silicate-polymerable organic compound epoxide compound) copolymer. In most products an excess amount of the aqueous alkali metal silicate may be used. The inorganic-organic plastic produced by the process of this invention has greatly improved flame resistance properties.\nThe polymerization of an alkali metal silicate with a polymerable unsaturated organic compound was illustrated in U.S. patent application Ser. No. 71,628, filed Sept. 11, 1970, by David H. Blount. The alkali metal silicate is oxidized by a peroxide initiator then polymerized with a polymerable organic compound. I have discovered that a stable emulsion of an aqueous alkali metal silicate and a polymerable unsaturated organic compound may be produced by adding up to 10% by weight, percentage based on weight of the aqueous alkali metal silicate solution, of a salt forming compound, and mixing with the mixture of the aqueous alkali metal silicate and polymerable compound. This stable emulsion greatly enhances the reaction between the alkali metal silicate and polymerable organic compound. Any suitable polymerable unsaturated organic compound may be used in this invention that can be polymerized in an aqueous alkali metal silicate solution in the presence of a peroxide initiator.\nThe emulsions of inorganic-organic plastics may be used as an adhesive on wood, paper, cement, plastics, ceramics, etc., as a coating agent on wood, cement, plastics, ceramics, etc., and may be dried or coagulated with a salt forming compound to produce a molding powder which may be molded by heat and pressure to produce useful objects such as knobs, handles, gears, pipes, toys, etc. The emulsion of inorganic-organic plastics may be further reacted with organic compounds such as polyisocyanates, isocyanates, epoxide compounds, substituted organic compounds, water-binding agents and many other compounds. The emulsion of inorganic-organic plastics may be used as a cavity filler, as putty, as a caulking compound, and in producing laminates.\nIt is accordingly, an object of my invention to provide novel inorganic-organic copolymers. A further object is to provide novel copolymers which may be used as an adhesive. A further object is to provide novel copolymers that will react with polyisocyanates to produce useful resinous and foam products. A further object is to provide a process for preparing novel inorganic-organic copolymers. Another object is to produce emulsions of inorganic-organic copolymers which may be used to produce concrete reinforced and reacted with inorganic-organic copolymers.\nThe inorganic-organic plastics may be produced by emulsifying and reacting the following components:\nComponent (a) an aqueous alkali metal silicate solution; PA0 Component (b) a polymerable unsaturated organic compound; PA0 Component (c) a salt forming compound; PA0 Component (d) an initiator PA0 Component (e) an epoxide compound PA0 p,p'-diphenylmethane diisocyanate PA0 phenylene diisocyanate PA0 chlorophenylene diisocyanate PA0 tolylene diisocyanate PA0 m-xylylene diisocyanate PA0 benzidine diisocyanate PA0 naphthylene diisocyanate PA0 tetramethylene diisocyanate PA0 pentamethylene diisocyanate PA0 hexamethylene diisocyanate PA0 decamethylene diisocyanate PA0 thiodipropyl diisocyanate"}
-{"text": "Electronic devices such as mobile phones, laptop computers, electronic notebooks, digital cameras, and video game instruments usually have an audio function. Such an electronic device having an audio function includes a housing, in which is placed a sound emitting part such as a speaker or a buzzer or a sound receiving part such as a microphone. The housing is provided with an opening at a position corresponding to that of the sound emitting part or the sound receiving part. Sounds are transmitted through the opening. A sound-transmitting membrane is attached to the housing so as to cover the opening of the housing in order to prevent foreign matters such as water droplets from entering the housing. The sound-transmitting membrane allows sound to pass through but prevents foreign matters from passing through. Known examples of the sound-transmitting membrane include porous plastic membranes such as porous polytetrafluoroethylene membranes and porous ultrahigh molecular weight polyethylene membranes (see Patent Literature 1).\nPatent Literature 2 describes a waterproof sound-transmitting member including a sound-transmitting membrane and a housing-side adhesive layer. The housing-side adhesive layer is a layer used to bond the waterproof sound-transmitting member to the housing and is laminated on the sound-transmitting membrane.\nPatent Literature 3 describes a waterproof sound-transmitting member including a sound-transmitting membrane, a support layer, and a housing-side adhesive layer. In the waterproof sound-transmitting member described in Patent Literature 3, the entire peripheral portion of the support layer extends outwardly beyond the housing-side adhesive layer."}
-{"text": "In the majority of television viewing environments today, when a user tunes to a specific channel, the graphical interface for that channel is defined primarily by the television manufacturer. For example, when a broadcaster displays a station identifying logo with their video transmission, they have no control over the positioning of their logo. Many times this logo is covered by a user selected menu option when overlaid on the television screen with the transmitted video.\nIn recent environments, graphical interfaces have provided a greater level of functionality with the advent of direct satellite system set top boxes (STB), such as DSS and Echostar. In addition to television displays, the user has channel banners, menus, electronic program guides and message alerts. In some circumstances, such as with inclement weather announcements and breaking news events, the user may customize the interface. However, the amount of functionality presented by these graphical interfaces is minimal since there is little interaction by a user, other than changing the channel or adjusting the volume.\nIn particular, these functions provided by an STB manufacturer occupy a portion of the display screen when selected by the user, such that the television broadcaster's logo or call letters are very likely to be covered up if displayed on the display screen at the same time a user is operating the electronic program guide, for example. The broadcasters have no control over the positioning of their station identifying information or other types of data they may wish to have in full view when displayed concurrently with user selected display options. With the DSS and Echostar systems, the additional user interfaces have been designed and implemented by STB manufacturers with little regard for the needs or concerns of the television broadcaster.\nTherefore, there is a direct conflict between the manufacturer's need to define the various functions provided by a direct satellite STB and the broadcaster's need to broadcast unobstructed graphical symbols and other information to viewers.\nFurthermore, this problem is not just limited to television broadcasts and receptions within the United States, but worldwide. This tension between the manufacturers and the broadcasters will increase as the boundaries of televisions and computing devices blur with the evolution of televisions adapting to the digital age and as data services become more readily available.\nThis convergence between televisions and computers can be illustrated by the efforts of Microsoft and Intel, where efforts have been made in developing strategies and approaches for integrating televisions into the personal computer (PC). For instance, Broadcast PC has been developed by Microsoft and Intel has developed Intercast, but they do not allow for the broadcaster's total control of the television environment. In both developments, a broadcaster is limited to controlling the look and feel of their services, but has no control of their primary product, which is the broadcast event.\nWindows CE, OS-9 and OpenTV are all operating systems that are attempting to deliver more sophisticated user interfaces to the manufacturer, but they too are disregarding the needs of the broadcaster. These operating systems still cover up the broadcast content with banners and menus that are genericized to the lowest common denominator of services and functions offered.\nTherefore, there is a need to permit broadcasters to have greater control in the television environment they transmit to their viewers. Specifically, there is a need for the broadcasters to broadcast unobstructed graphical symbols and other information to viewers."}
-{"text": "As the supply of low sulfur, low nitrogen crudes decrease, refineries are processing crudes with greater sulfur and nitrogen contents at the same time that environmental regulations are mandating lower levels of these heteroatoms in products. Consequently, a need exists for increasingly efficient desulfurization and denitrogenation catalysts.\nIn one approach, a family of compounds, related to hydrotalcites, e.g., ammonium nickel molybdates, has been prepared. Whereas X-ray diffraction analysis has shown that hydrotalcites are composed of layered phases with positively charged sheets and exchangeable anions located in the galleries between the sheets, the related ammonium nickel molybdate phase has molybdate anions in interlayer galleries bonded to nickel oxyhydroxide sheets. See, for example, Levin, D., Soled, S. L., and Ying, J. Y., Crystal Structure of an Ammonium Nickel Molybdate prepared by Chemical Precipitation, Inorganic Chemistry, Vol. 35, No. 14, p. 4191-4197 (1996). The preparation of such materials also has been reported by Teichner and Astier, Appl. Catal. 72, 321-29 (1991); Ann. Chim. Fr. 12, 337-43 (1987), and C. R. Acad. Sci. 304 (II), #11, 563-6 (1987) and Mazzocchia, Solid State Ionics, 63-65 (1993) 731-35.\nNow, when molybdenum is partially substituted for by tungsten, an amorphous phase is produced which upon decomposition and, preferably, sulfidation, provides enhanced hydrodenitrogenation (HDN) catalyst activity relative to the unsubstituted (Ni--Mo) phase."}
-{"text": "1. Field of the Invention\nThe present invention relates to a microprocessor that can prevent illegal alternation of execution codes and processing target data under a multi-task program execution environment.\n2. Description of the Background Art\nIn recent years, the performance of a microprocessor has improved considerably such that the microprocessor is capable of realizing reproduction and editing of video images and audio sounds, in addition to the conventional functions such as computations and graphics. By implementing such a microprocessor in a system designed for end-user (which will be referred to as PC hereafter), the users can enjoy various video images and audio sounds on monitors. Also, by combing the function for reproducing video images and audio sounds with the computational power of the PC, the applicability to games or the like can be improved. Such a microprocessor is not designed for any specific hardware and can be implemented in a variety of hardwares so that there is an advantage that the users who already possess PCs can enjoy reproduction and editing of video images and audio sounds inexpensively by simply changing a microprocessor for executing programs.\nIn the case of handling video images and audio sounds on PCs, there arises a problem of a protection of the copyright of original images or music. In the MD or digital video playback devices, unlimited copies can be prevented by implementing a mechanism for preventing the illegal copying in these devices in advance. It is rather rare to attempt the illegal copying by disassembling and altering these devices, and even if such devices are made, there is a worldwide trend for prohibiting the manufacturing and sales of devices altered for the purpose of illegal copying by laws. Consequently, damages due to the hardware based illegal copying are not very serious.\nHowever, image data and music data are actually processed on the PC by the software rather than the hardware, and the end-user can freely alter the software on the PC. Namely, if the user has some level of knowledge, it is quite feasible to carry out the illegal copying by analyzing programs and rewriting the executable software. In addition, there is a problem that the software for illegal copying so produced can be spread very quickly through media such as networks, unlike the hardware.\nIn order to resolve these problems, conventionally a PC software to be used for reproducing copyright protected contents such as commercial films or music has employed a technique for preventing analysis and alternation by encrypting the software. This technique is known as a tamper resistant software (see David Aucsmith et al., \u201cTamper Resistant Software: An Implementation\u201d, Proceedings of the 1996 Intel Software Developer's Conference).\nThe tamper resistant software technique is also effective in preventing illegal copying of valuable information including not only video and audio data but also text and know-how that is to be provided to a user through the PC, and protecting know-how contained in the PC software itself from analysis.\nHowever, the tamper resistant software technique is a technique which makes analysis using tools such as de-assembler or debugger difficult by encrypting a portion of the program that requires protection before the execution of the program starts, decrypting that portion immediately before executing that portion and encrypting that portion again immediately after the execution of that portion is completed. Consequently, as along as the program is executable by a processor, it is always possible to analyze the program by carrying out the analysis step by step starting from the start of the program.\nThis fact has been an obstacle for a copyright owner to provide copyright protected contents to a system for reproducing video and audio data using the PC.\nThe other tamper resistant software applications are also vulnerable in this regard, and this fact has been an obstacle to a sophisticated information server through the PC and an application of a program containing know-how of an enterprise or individual to the PC.\nThese are problems that equally apply to the software protection in general, but in addition, the PC is an open platform so that there is also a problem of an attack by altering the operating system (OS) which is intended to be a basis of the system's software configuration. Namely, a skilled and malicious user can alter the OS of his own PC to invalidate or analyze the copyright protection mechanisms incorporated in application programs by utilizing privileges given to the OS.\nThe current OS realizes the management of resources under the control of the computer and the arbitration of their uses by utilizing a privileged operation function with respect to a memory and an execution control function provided in CPU. Targets of the management include the conventional targets such as devices, CPU and memory resources, as well as QoS (Quality of Service) at network or application level. Nevertheless, the basics of the resource management are still allocations of resources necessary for the execution of a program. Namely, an allocation of a CPU time to the execution of that program and an allocation of a memory space necessary for the execution are the besics of the resource management. The control of the other devices, network and application QoS is realized by controlling the execution of a program that makes accesses to these resources (by allocating a CPU time and a memory space).\nThe OS has privileges for carrying out the CPU time allocation and the memory space allocation. Namely, the OS has a privilege for interrupting and restarting an application program at arbitrary timing and a privilege to move a content of a memory space allocated to an application program to a memory of a different hierarchical level at arbitrary timing, in order to carry out the CPU time allocation. The latter privilege is also used for the purpose of providing a flat memory space to the application by concealing (normally) hierarchical memory systems with different access speeds and capacities from the application.\nUsing these two privileges, the OS can interrupt an execution state of the application and take a snap shot of it at arbitrary timing, and restart it after making a copy of it or rewriting it. This function can also be used as a tool for analyzing secrets hidden in the application.\nIn order to prevent an analysis of the application on a computer, there are several known techniques for encrypting programs or data (Hampson, U.S. Pat. No. 4,847,902; Hartman, U.S. Pat. No. 5,224,166; Davis, U.S. Pat. No. 5,806,706; Takahashi et al., U.S. Pat. No. 5,825,878; Buer et al., U.S. Pat. No. 6,003,117; Japanese Patent Application Laid Open No. 11-282667 (1999), for example). However, these known techniques do not account for the protection of the program operation and the data secrecy from the above described privileged operations of the OS.\nThe conventional technique based on the x86 architecture of Intel Corporation (Hartman, U.S. Pat. No. 5,224,166) is a technique for storing the execution codes and data by encrypting them by using a prescribed encryption key Kx. The encryption key Kx is given in a form of EKr[Kx] which is encrypted by using a public key Kp corresponding to a secret key Ks embedded in a processor. Consequently, only the processor that knows Ks can decrypt the encrypted execution codes on a memory. The encryption key Kx is stored in a register inside the processor called a segment register.\nUsing this mechanism, it is possible to protect the secrecy of the program codes from the user to some extent by encrypting the codes. Also, it becomes cryptographically difficult for a person who does not know the encryption key Kx of the codes to alter the codes according to his intention or newly produce codes that are executable when decrypted by using the encryption key Kx.\nHowever, the system employing this technique has a drawback in that the analysis of the program becomes possible by utilizing a privilege of the OS called a context switching, without decrypting the encrypted execution codes.\nMore specifically, when the execution of the program is stopped by the interruption or when the program voluntarily calls up a software interruption command due to the system call up, the OS carries out the context switching for the purpose of the execution of the other program. The context switching is an operation to store an execution state (which will be referred to as a context information hereafter) of the program indicating a set of register values at that point into a memory, and restoring the context information of another program stored in the memory in advance into the registers.\nFIG. 15 shows the conventional context storing format used in the x86 processor. All the contents of the registers used by the application are contained here. The context information of the interrupted program is restored into the registers when the program is restarted. The context switching is an indispensable function in order to operate a plurality of programs in parallel. In the conventional technique, the OS can read the register values at a time of the context switching, so that it is possible to guess most of the operations made by the programs if not all, according to how the execution state of that program has changed.\nIn addition, by controlling a timing at which the exception occurs by setting of a timer or the like, it is possible to carry out this processing at arbitrary execution point of the program. Apart from the interruption of the execution and the analysis, it is also possible to rewrite the register information by malicious intention. The rewriting of the registers can not only change the operation of the program but also make the program analysis easier. The OS can store arbitrary state of the application so that it is possible to analyze the operation of the program by rewriting the register values and operating the program repeatedly. In addition to the above described functions, the processor has a debugging support function such as a stepwise execution, and there has been a problem that the OS can analyze the application by utilizing all these functions.\nAs far as data are concerned, U.S. Pat. No. 5,224,166 asserts that the program can access the encrypted data only by the program execution using the encrypted code segment. Here, there is a problem that the encrypted data can be freely read by the encrypted program by using arbitrary key, regardless of the encryption key by which the program is encrypted, even when there are programs encrypted by using mutually different encryption keys. This conventional technique does not account for the case where the OS and the application have their own secrets independently and the secret of the application is to be protected from the OS or a plurality of program providers have their own secrets separately.\nOf course, it is possible to separate memory spaces among the applications and to prohibit accesses to a system memory by the applications by the protection function provided in the virtual memory mechanism even in the existing processor. However, as long as the virtual memory mechanism is under the management of the OS, the protection of the secret of the application cannot rely on the function under the management of the OS. This is because the OS can access data by ignoring the protection mechanism, and this privilege is indispensable in providing the virtual memory function as described above.\nAs another conventional technique, Japanese Patent Application Laid Open No. 11-282667 (1999) discloses a technique of a secret memory provided inside the CPU in order to store the secret information of the application. In this technique, a prescribed reference value is required in order to access data in the secret memory. However, this reference fails to disclose how to protect the reference value for obtaining the access right with respect to the secret data from a plurality of programs operating in the same CPU, especially the OS.\nAlso, in U.S. Pat. No. 5,123,045, Ostrovsky et al. disclose a system that presupposes the use of sub-processors having unique secret keys corresponding to the applications, in which the operation of the program cannot be guessed from the access pattern by which these sub-processors are accessing programs placed on a main memory. This is based on a mechanism for carrying out random memory accesses by converting the instruction system for carrying out operations with respect to the memory into another instruction system different from that.\nHowever, this technique requires different sub-processors for different applications so that it requires a high cost, and the implementation and fast realization of the compiler and processor hardware for processing such instruction system are expected to be very difficult as they are quite different from those of the currently used processors. Besides that, in this type of processor, it becomes difficult to comprehend correspondences among the data contents and the operations even when the data and the operations of the actually operated codes are observed and traced so that the debugging of the program becomes very difficult, and therefore this technique has many practical problems, compared with the other conventional techniques described above in which the program codes and the data are simply encrypted, such as those of U.S. Pat. No. 5,224,166 and Japanese Patent Application Laid Open No. 11-282667."}
-{"text": "1. Field of Invention\nThis invention relates to the technical field of mimeograph, and more particularly to a method of controlling a linkage drive section in a stencil printing machine to set a printing drum in place in which the position of a print image can be adjusted in a horizontal direction with ease. In this specification, the term \"horizontal direction\" is intended to mean \"a direction which is in parallel with the central axis of the printing drum\"; and the term \"horizontal position\" is intended to mean \"a position along the direction thus defined\".\n2. Description of Related Art\nIn a conventional rotary stencil printing machine with a printing drum, the printing drum is set at a predetermined position in the direction of its axis in the printing machine frame and rotated around its central axis. In the machine, adjustment of the horizontal position of a print image can be achieved by suitably shifting the position of the mimeographic stencil paper on the printing drum in the direction of the axis of the printing drum before it is wound on the printing drum. Alternatively, by moving the sheet supplying device in the direction of the axis of the printing drum, the horizontal position of the print image can be adjusted at all times, not only in the case where the stencil paper is not wound on the printing drum yet but also in the case where it is wound on it.\nIn moving a mimeographic stencil paper or a printing sheet, which is rectangular and has two short sides extended in the direction of movement and two long sides which are perpendicular to the short sides, over the printing drum along the axis of the latter, the stencil paper or printing sheet must be parallel-moved in the direction of the short sides with high accuracy. In order to make this movement delicate and accurate thereby to permit the fine adjustment in horizontal position of the print image, means for horizontally moving a stencil paper feeding device adapted to feed a stencil paper to the printing drum and a sheet supplying device adapted to supply printing sheets to the printing drum should have: guide means which use a pair of parallel guide rails laid in the lateral direction of the stencil paper or printing sheet to guide the stencil paper feeding device or the sheet supplying device accurately parallel at least two positions spaced from each other in the longitudinal direction of the stencil paper or printing sheet; and fine drive means for giving feeding actions to the stencil paper feeding device or the sheet supplying device along the guide means for the same distance at the same speed.\nAs is apparent from the above description, in the case where, with respect to the printing drum which has been fixed in the axis direction, the stencil paper or printing sheet is moved along the axis of the printing drum to adjust the horizontal position of the print image, it is necessary to provide considerably intricate means."}
-{"text": "It is well recognized that combating oil spills in coastal waters is of great communal interest. This has led to continuous efforts to improve oil spill mitigation strategies. The preferred spill mitigation strategy often uses oil booms to contain and concentrate floating oil, prior to oil skimmer recovery. Despite improvements over the decades, limitations of tow speed and operation under real world conditions remain serious.\nConventional Oil Boom Leakage\nA well-designed oil boom should be flexible to conform to wave motions, yet sufficiently rigid to retain as much oil as possible. Designs range from small, lightweight booms for manual harbor deployment, to large, robust booms for open sea use that need a crane and sizeable vessels to handle and deploy.\nImportant oil boom failure modes include overtopping where oil passes over the boom, and leakage where oil passes under the boom, surfacing as downstream oil patches. Boom overtopping mechanisms include overfilling, wave splashover, and boom-diving due to high towing speed. Boom leakage includes overfilling, frontal slick droplet injection (entrainment) from the confined pool, and boom-surfing due to opposing currents and/or winds blowing in the tow direction.\nOil pool thickness at the boom increases as more oil is collected and/or as the boom is towed faster. Overfilling often occurs with waves and current-induced boom diving and splashover. Leakage occurs when oil droplets are injected sufficiently deep to pass under the oil boom and is highly sensitive to towing speed. Once the towing velocity exceeds the boom's critical towing speed, the frontal wave \u2018breaks\u2019 on a large scale, increasing oil entrainment dramatically. Instability development and thus droplet injection increase strongly with towing speed.\nLeakage also occurs for acute angle flows at boom-segment junctions from vortices formation. These vortices can inject oil into the water that then underpasses the boom. Acute angle boom flows occur on the leading boom segments that steer oil towards the apex.\nBubble Oil Boom Background\nIn 1971, the U.S. Coast Guard tested a pneumatic boom for oil spill control, otherwise known as a bubble oil boom (BOB) (US Coast Guard. Heavy-duty oil containment systems: pneumatic barrier system. Report 714102/A/094, US Coast Guard Office of Research and Development, Contract DOT-CG-00-490-A to Oil Containment Division, Wilson Industries, 1971). A pneumatic boom generally is formed from a long, submerged air pipe with a series of holes along its length, typically at the upper generatrix. The curtain of bubbles rise in a sheet that drives an upwelling flow, which at the water surface is converted (by continuity) into an outwelling, which is the oil-blocking barrier. The Coast Guard study concluded that BOBs only were useful for low currents such as harbors, where current commercial BOBs are found. A number of other attempts at using BOBs have been made, such as in U.S. Pat. Nos. 3,491,023, 3,744,254, and more recently U.S. Patent Publication No. 2011/0303613. However, to date there are no large-scale commercial BOB systems available outside harbors due a number of drawbacks with previous designs."}
-{"text": "This invention relates to gas detecting apparatus, and more particularly, to natural gas detecting apparatus which is able to discriminate between pipeline natural gas, non-pipeline sources of methane, propane, and gasoline vapors.\nUtilities which distribute natural gas require reliable gas-leak detectors for use in maintenance of gas supply lines. Existing natural gas detectors are either costly, sensitive and non-selective or low cost, insensitive and non-selective. Non-selective gas detectors respond to any combustible gas. Selective gas detectors are specific to hydrocarbon gases. The two presently most used gas detectors are based on hydrogen flame ionization and on hot wire catalysis. These gas detectors cannot distinguish among different types of hydrocarbons. However, it is necessary to distinguish among different types of hydrocarbons in order to distinguish a pipeline gas from gasoline vapors or sewer or swamp gas and so reduce leak surveyor time wasted on false alarms. Ethane content, if measurable, provides a good means to discriminate between pipeline gas and interfering gasoline vapors and sewer or swamp gases because the later contain practically no ethane, while pipeline gas does, in varying degrees. Gasoline vapors and propane (LP gas) can also generate a false alarm with conventional instruments. However, their infrared absorption is shifted relative to that of methane, as will be described later, as is the basis for this invention to eliminate such false alarms.\nIn U.S. Pat. No. 4,507,558, there is disclosed a selective detector for natural gas which discriminates between low concentrations of natural gas and other methane sources by measuring the characteristics of the methane/ethane ratio of natural gas as well as by using a combustible gas sensor. The operation of this detector is based on infrared light absorption of methane and ethane in combination with another non-specific combustible gas detector whereby the detector has the ability to detect nonspecifically, the presence of a combustible gas, and to define the nature of the combustible gas. Thus, this natural gas detector utilizes two types of detection including nondispersive infrared detectors and a non-specific combustible detectors such as hot-wire catalytic combustible detector. The detector determines concentration of both methane or ethane irrespective of the concentration of the other gas by using absorption cells placed in front of the detectors. The detector includes a light emitting diode which issues light centered around 3.32 microns and a reference light source which emits light at a wavelength outside of this band. Although this arrangement permits distinguishing among different types of hydrocarbons, the requirement for a hotwire catalytic combustible detector adds cost and complexity to the device and increases power consumption.\nIt would be desirable to have a natural gas detector which can distinguish among different types of hydrocarbons, and which provides information to the user on the amount and type of combustible gases in the environment."}
-{"text": "The present invention relates to a magneto-rheological (xe2x80x9cMRxe2x80x9d) fluid damper, and more particularly, to a linearly-acting MR fluid damper suitable for vibration damping in a vehicle suspension system.\nMR fluids are materials that respond to an applied magnetic field with a change in Theological behavior (i.e., change in formation and material flow characteristics). The flow characteristics of these non-Newtonian MR fluids change several orders of magnitude within milliseconds when subjected to a suitable magnetic field. In particular, magnetic particles noncolloidally suspended in fluid align in chain-like structures parallel to the applied magnetic field, changing the shear stress on adjacent shear surfaces.\nDevices such as controllable dampers benefit from the controllable shear stress of MR fluid. For example, linearly-acting MR fluid dampers are used in vehicle suspension systems as vibration dampers. At low levels of vehicle vibration, the MR fluid damper lightly damps the vibration, providing a more comfortable ride, by applying a low magnetic field or no magnetic field at all to the MR fluid. At high levels of vehicle vibration, the amount of damping can be selectively increased by applying a stronger magnetic field. The controllable damper lends itself to integration in vehicle suspension systems that respond to vehicle load, road surface condition, and driver preference by adjusting the suspension performance.\nIn some applications, linearly-acting MR fluid dampers use a piston assembly that moves within a damper body tube having a cylindrical reservoir that separates a volume of MR fluid into a compression chamber and an extension chamber. The piston assembly has a piston core positioned within a flux ring to form an annular flow gap therebetween. Relative motion between the damper body tube and the piston assembly is dampened by a flow of the MR fluid through the flow gap from one chamber to another caused by the relative motion.\nAlignment of the flux ring is critical for optimum performance. Ideally, the piston assembly should move freely within the reservoir in the damper body tube without friction or binding. In addition, the radial width and concentricity of the annular flow passage must be precisely set and maintained along the axial length of the passage throughout the operation to ensure optimum, predictable control of the damping. Consequently, the flux ring must be correctly aligned with the piston core.\nAttachment elements have been suggested to provide flux ring alignment with nonmagnetic bridge elements. In particular, perforated end plates are aligned above and below the flux ring and piston core. These attachment elements have several potential problems. First, the attachment elements increase the length of the piston assembly. Consequently, less travel distance is available for the piston to move within the cylindrical reservoir of the damper body tube. Second, the attachment elements require tight manufacturing tolerances in order to correctly align the flux ring to the piston core. Third, such attachment elements often include tabs or other projections that increase the drag as the piston moves, which may be undesirable. Fourth, the attachment elements have numerous components and require manufacturing operations such as spot welding. Therefore, such attachment elements are costly to manufacture and time consuming to assemble.\nConsequently, there is a need for an improved piston assembly suitable for use in a magneto-rheological (MR) fluid damper.\nThe present invention addresses the above need by providing an improved piston assembly for a linearly-acting MR fluid damper. The piston assembly of the present invention confines a flux ring within the functional length of the piston assembly without significantly restricting fluid flow, thereby providing optimum performance with minimal piston length. Further, the part count of the piston assembly is reduced; and the piston assembly is easier to assemble in a desired alignment. Thus, the piston assembly of the present invention is of a simpler construction than known damper pistons that can be manufactured for less cost.\nAccording to the principles of the present invention and in accordance with the described embodiment, the present invention provides a piston assembly for use with an MR fluid damper. The piston assembly has a flux ring positioned in a desired alignment with a piston core to form an annular flow gap between the flux ring and the piston core. The piston core is secured to the flux ring in the desired alignment by a plurality of projections extending across the flow gap between an inner surface of the flux ring and an outer surface of the piston core. Thus, the flux ring is secured on the piston ring without using expensive, high precision attachment components; and the piston assembly is able to utilize its full length, thereby providing optimum performance with a minimum of length piston.\nIn one aspect of the present invention, the projections are molded through attachment passages in the flux ring. In a further aspect of the invention, the attachment passages are holes intersecting the inner and outer surfaces of the flux ring.\nIn another embodiment of the invention, a method is provided for making a piston assembly for use with an MR fluid damper. The method comprises first, fixing a piston core and a flux ring in a desired alignment forming a flow gap therebetween; and then, forming a plurality of projections in the flow gap between the flux ring and the piston core to secure the flux ring on the piston core in a desired alignment."}
-{"text": "Soft drink bottles are ordinarily packaged by bottlers in cases holding several plastic or glass bottles for shipment to retailers or for storage. The cases are customarily stacked on top of the bottles in the lower tray for storage in warehouses. The trays must therefore be particularly stable in order to remain standing in storage. The bottles are generally stored in plastic low depth cases in which the side walls of the cases are lower than the height of the bottles and in which the bottles support the weight of the cases stacked on top of the bottles. When cases of these are stacked for display, the only bottles which can be removed from the stack are those stored in the top case.\nCurrently in order for a bottle to be displayed, it must be removed from the tray and placed in a display rack. If the trays are stacked on the bottles for display, the bottles can only be removed from the top tray. The top tray is removed from the stack when empty to provide access to the bottles on the next tray. Trays stacked on the tops of the bottles in the lower trays are unstable, particularly when bottles are removed from the lower trays while in storage."}
-{"text": "1. Technical Field\nThe present invention relates to an image forming apparatus, in which a toner image formed on a transfer member, based on image data, is transferred and fused on a recording medium to form an image.\n2. Description of the Background Art\nIn electrophotographic image forming apparatuses, toner images are formed on a transfer belt using image data, and then transferred and fused on a sheet to form images. Such image forming apparatuses can form images using colored toner such as yellow, magenta, cyan, and black toner, and also clear toner. For example, an image forming apparatus can form images having a watermark on the top layer of the images by superimposing a toner image of clear toner over the toner images of colored toner, and transferring and fusing each of toner images on a sheet.\nCompared to color image formation using only the colored toner, when an image is formed using a combination of the colored toner and the clear toner, the total amount of toner to form the image becomes great. However, in image forming apparatuses, if the total amount of toner transferred on the sheet becomes great, heat amount used for plasticizing or melting toner for the normal transfer process may not be enough for fusing a toner image on the sheet, and resultantly a transfer failure may occur.\nIn light of such problem, JP-2009-63744-A discloses an image forming apparatus using an image forming method in which an upper limit (or control value) is set for the total combined amount of colored toner and clear toner, and when the total amount of toner exceeds the upper limit, the density or concentration of colored toner is adjusted. In such a method, when the total amount of toner exceeds the upper limit, the density or concentration of clear toner is fixed to a given value, and the density or concentration of colored toner is decreased to limit the total amount of toner at the upper limit. With such a configuration, an image can be formed without affecting the gloss provided by the clear toner.\nHowever, such toner-amount reduction of the colored toner amount has the undesirable effect of decreasing image density of the resultant output image."}
-{"text": "In general, current cholangiographic agents used as biliary X-ray contrast agents are now considered too toxic to meet risk/benefit criteria for routine use as contrast agents to visualize the biliary system. Because of the relative toxicity of these agents, these agents are being replaced in use by alternative imaging methodologies, such as ultrasound and magnetic resonance imaging (MRI). On the other hand, the resolution achievable with X-ray radiology techniques is superb. Accordingly, biliary contrast agents can still be quite useful.\nA second issue with regard to these agents is the low rate of choleresis associated with biliary excretion of the agent. The low rate of choleresis increases the biliary concentration and thus increases the toxicity.\nCurrent biliary X-ray contrast agents have amphophilic properties. These agents have a hydrophilic head portion; that is, the --CO.sub.2 H terminus, and a hydrophobic tail; that is, triiodo aniline moiety with an unsubstituted 5-position. Because of this property, they are excreted nearly exclusively by the biliary route. Typical of most amphilphilic molecules, these biliary contrast agents tend to self associate. This property is reflected by the low critical micellar concentration and low osmolalities of their aqueous salt solution.\nIt is therefore desirable to produce a less toxic and more choleretic contrast agent."}
-{"text": "All musical instruments, especially multi-stringed ones such as the piano and harpsichord, require periodic tuning to insure that they consistently reproduce the proper pitches when played.\nAt present this periodic tuning is accomplished primarily by listening for the beat signal between the note and an audible reference frequency. This is inherently an inaccurate method, and in the case of multi-stringed instruments such as pianos and harpsichords it requires skill far beyond that of the average musician.\nThere exist today two types of electronic aids which enable the person of no special skill to successfully tune a multi-stringed instrument. In the first of these aids the audible note from the instrument is converted into pulsations of light which are used to illuminate a rotating strobe disc. Deviation of the frequency of the note from the proper value is displayed as apparent movement of a pattern of light and dark spots on the disc.\nThe second method employs an electronic frequency counter to measure the frequency of the note and provide a digital display of the frequency in hertz. The user must then compare this frequency to the desired frequency and adjust the instrument accordingly.\nBoth these methods suffer from inherent disadvantages which have prevented their widespread use. The strobe disc method employs an electric motor and is therefore bulky as well as being relatively inaccurate. The frequency counter method has a slow response time, limiting its usefulness to the middle and upper octaves, in addition to which the digital readout is difficult for an untrained person to interpret.\nThe invention herein described seeks to overcome these disadvantages, thus making it possible for an untrained person to quickly tune any note in the musical spectrum with an accuracy exceeding that of the best professional ear."}
-{"text": "1. Field of the Invention\nMy invention relates generally to an adjustable box wrench and particularly to a box wrench that positively grips the workpiece in response to torque applied to a cam lever separately articulated from the wrench body.\n2. Discussion of the Related Art\nThe box-end or box wrench is known in the art for turning polygonal nuts and bolts located in tight and inaccessible locations. The box wrench usually includes a thin-walled polygonally-fluted box structure on one or both ends. Because this thin-walled box structure is sized to fit snugly over the generally polygonal nut or bolt workpiece, the workpiece can be engaged for torquing even where there is only a few millimeters of clearance.\nAdjustable wrenches, including adjustable box-end wrenches, are also known in the art. For instance, in U.S. Pat. No. 4,325,275, David S. Colvin discloses an adjustable open-end and box-end wrench that uses a pair of spaced pins and skewed slots to allow adjustment of a member that cooperates with a box-end to re-size the box structure for a range of workpiece dimensions. Although Colvin's box-end wrench can be used to engage a workpiece in tight spaces, he neither considers nor suggests \"gripping\" means for positively engaging a deformed workpiece and instead relies on the matching polygonal geometry of box structure and workpiece to transfer torque therebetween. Similarly, in U.S. Pat. No. 3,858,465, Erik Lind discloses an adjustable wrench employing an axially displaceable and lockable external sleeve to adjust the geometry of the box structure but neither considers nor suggests means for positively gripping a deformed workpiece. In U.S. Pat. No. 3,363,490, K. Maichen discloses a double-ended simultaneously-adjustable wrench that operates by manually turning an eccentric pin to move two members cooperating to form surfaces for receiving and torquing polygonal workpieces. Maichen also neither considers nor suggests means for positively gripping a deformed workpiece. Swiss Patent No. 386 948 issued to August Samuel Aegerter and U. K. Patent 251,544 issued to Andrew Arbuckle both disclose adjustable box-end wrenches that rely on the manual rotation of a cam to urge a sliding adjusting rod into a position that fixes the effective engagement dimensions of a box structure. Neither Aegerter nor Arbuckle consider or suggest means for positively gripping a deformed workpiece.\nThe common problem of transferring effective torque to deformed nuts and bolts is well-known in the art. When a nut is \"frozen\" onto a threaded bolt, the torque applied in attempting to remove it may deform the polygonal geometry of the outer nut surface to such an extent that a common box-end wrench (even an adjustable one) no longer properly engages the deformed surface to transfer torque. Rounded comers merely slip within the engaging box structure when torque is applied, accomplishing nothing.\nGripping pliers and wrenches are known in the art for torquing deformed workpieces. For instance, in U.S. Pat. No. 3,611,843, Joachim E. Engel discloses an adjustable socket wrench that has a gripping handle and a pair of relatively movable jaws, one of which is secured to the gripping handle and the other of which is cammed into engagement with its neighbor. Engel's handle includes a coarse adjusting member coupled to an axially movable jaw to permit adjustment of the socket dimensions by turning a threaded member. Engel's wrench increases the gripping force in the jaws responsive to an increase in the pivotal rotation of the jaws relative to the gripping handle arising from torque applied to the handle but cannot grip deformed workpieces in tight spaces. Similarly, in U.S. Pat. No. 4,174,646, Simon Cotler discloses a universal tool with gripping action and replaceable jaws that has a body with an opening adapted to interchangeably receive cassette-type work elements such as a box-end wrench element or the like. Cotler uses a cam-locking lever integrated into the universal tool to apply a predetermined gripping force on the workpiece. Disadvantageously, his universal tool employs a bulky structure to receive the cassette-type working element and is not suited for gripping polygonal workpieces in tight spaces. Moreover, contrary to Engel's advantageous feature of increasing gripping force proportionately with applied torque, Cotler's cam-locking scheme provides an unvarying gripping force that may be insufficient to retain the workpiece at high torque levels. Finally, in U.S. Pat. No. 2,486,523, P. E. Deschenes discloses a similar cam-locking adjustable gripping wrench for use with bottle caps.\nThus, although gripping wrenches are known in the art, including cam-operated and cam-locking gripping pliers and wrenches, no gripping wrench suitable for use in the tight spaces serviced by box-end wrenches has been known in the art until now. A wrench adapted to gripping and torquing polygonal workpieces in tight spaces could satisfy a clearly-felt need in the art. The related unresolved problems and deficiencies are clearly felt in the art and are solved by my invention in the manner described below."}
-{"text": "It is well known that fluorosilicones are used as lubricant oils and base fluids for greases. They also show good rheological properties and resistance to degradation and oxidation. However, this lubricant behavior, which is satisfactory at relatively low temperatures (e.g., up to 100.degree. C.) tends to progressively worsen at higher temperatures (e.g., from 100.degree. to 200.degree. C. and higher). The prior art has taught the use of various additives, such as antioxidants and/or antiwear agents, in an attempt to extend the practical application temperature range of fluorosilicones. But, e.g., Braun et al., in an article entitled \"Silicone Lubrication of Porous Bronze Bearings\" (Lub. Eng., 32, 176-182, 1975) found some antioxidants to be ineffective. Further, some fatty ester antiwear agents were found to give similar unsatisfactory results.\nKim et al. (U.S. Pat. No. 3,629,115) suggested the use of fluorinated phenylphosphine as both an antioxidant and an antiwear additive in fluorosilicones. However, the antiwear properties of these compositions are not satisfactory in practice.\nKobzova et al., in an article entitled \"Efficiency of Antioxidants in Phthalocyanine Greases\" (Khim. i Tekhn. Topliv i Masel, No. 10, pp 59-61, 1971) showed that conventional antioxidants, such as N-phenyl-.alpha.-naphthylamine, are nearly ineffective in fluorosilicone-based greases.\nTherefore, from an application standpoint, there is still a need for agents which, when added to fluorosilicones, result in effective lubricants for use in air at temperatures higher than 100.degree. C., particularly at 125.degree. C. to 200.degree. C., and higher."}
-{"text": "1. Field of the Invention\nThe invention relates to web applications over the internet, and in particular, to a disruption-free web application execution mechanism.\n2. Description of the Related Art\nA web application is a program providing services through the hyper text transport protocol (HTTP), allowing a client to access services on a remote server over the internet. One field utilizing web applications is mobile communications. For mobile communications, temporarily disconnecting and reconnecting to the internet from a remote service frequently occurs, because a client is usually connected through variable environments where transmission signal quality may be uncertain. Thus, various offline execution methods are proposed to ensure that a remote service can be provided without disruption.\nConventional offline services are mostly implemented by specific modified programs provided by the remote server, wherein offline functions are defined should disconnection occurs, such as defined offline functions provided by Google Gears and Adobe AIR. The disadvantage of the approach is that since the offline functions require a specific modified program to be executed, flexibility for clients is decreased and clients must install the specific modified program for compatibility. Thus, a more compatible mechanism is desirable."}
-{"text": "1. Field of the Invention\nThe present invention relates to a document illuminating system for use in an image reader such as a digital copying machine or an image scanner to illuminate an object such as a document as well as to an image reader incorporating such a document illuminating system. It is applicable to a digital copying machine, a facsimile machine, a printer, and a complex machine of these machines.\n2. Description of the Related Art\nIn recent years, LEDs (Light Emitting Diodes) have been well developed and highly bright LEDs are now available, LEDs offer such advantages as high longevity, high efficiency, high resistance to impact, and a mono-color light emission and will be expectedly adapted to various types of illumination including an illumination system of an image reader as a digital copying machine or an image scanner. Further, emission spectrum of a white LED covers a visible range of wavelength bands so that it is adoptable for a document illuminating system of a color image reader.\nFor this reason, various kinds of document illuminating systems using LEDs have been proposed. There is a well known technique to properly converge light beams with diffusion property from LEDs by a light guide plate or a reflective mirror in order to efficiently illuminate a document read area.\nFor example, Japanese Laid-open Patent Publication No. 2006-42016 (Reference 1) and No. 2006-295810 (Reference 2) disclose a document illumination system and an image reader to improve illumination performance of an illumination system by use of a reflective mirror. In both of the References, the position of the illumination system is diagonal relative to a document, the illumination system is configured to directly illuminate a document with a light component of a light beam from the LEDs at a small divergent angle to a direction orthogonal to a light emission face and to guide a light component at a large divergent angle to the document by use of a reflective mirror. Reference 1 exemplifies a structure with a reflective mirror including a curved face while Reference 2 exemplifies a structure with a reflective mirror including a planar face.\nHowever, there is a problem in the References 1 and 2 that the diagonal position of the illumination system to the document leads to an increase in thickness and size of the illumination system in an orthogonal direction to the document. In order to decrease the size of the reader, a distance from a light source to the document need to be decreased, which makes it impossible to illuminate the document with an even luminance distribution. Furthermore, this illumination system is configured to illuminate the document from a single direction so that dark lines may occur in an image in reading an uneven document.\nIn addition, Reference 2 discloses an LED structure having a special covering element to cover an emission face, however, such a covering element faces a problem of high manufacture cost since general low-cost LEDs are not suitable for such a structure.\nIn order to deal with occurrence of dark lines in an image, Japanese Laid-open Patent Publication No. 2008-67276 (Reference 3) discloses a structure to illuminate an object with light from two illumination elements which oppose to each other (hereinafter, sometimes referred to as opposite illumination). This can prevent dark lines in an image of an uneven document illuminated with light from a single direction.\nHowever, the two opposite illumination elements are each provided with light sources, which also results in increases in the size and manufacture costs of the illuminating system. From the sub-scan direction cross sectional (orthogonal to the document and including sub scan direction) drawing of the Reference 3, it is seen that the circuit board of light emitting elements extends diagonally downwards. Therefore, the document illuminating system of Reference 3 is increased in size in a height direction thereof. Moreover, since a normal line of the light emitting element directs to a direction of the document read area, it is difficult to secure a sufficient optical path length, leading to enlargement and uneven luminance of the illuminating system.\nIn view of solving the above problems, there is a demand for development of a document illuminating system in small thickness which can oppositely illuminates documents without occurrence of dark lines in an image.\nMoreover, along with an increasing demand for a higher-speed, higher image quality image reader, it is required to develop a document illuminating system configured to guide light onto the light-receiving face of the image sensor in even luminance distribution in the sub scan direction for the following reasons.\nA read line width of a CCD on a document corresponds to a size of the document in a main scan direction, and a position thereof in the sub scan direction may be shifted from a referential sub scan position due to an error in adjustment of such an optical element as a reflective mirror and an imaging lens or a movement of respective elements in reading a document from one end to the other in the main scan direction. Generally, the shift of a read line width may be about 0.2 to 1.0 mm in the sub scan direction relative to the optical axis of a read system.\nBecause of this, without illuminating a document with light amounts in even distribution, there will be an area insufficiently illuminated or an area over-illuminated in the document and a defective image (decrease in read accuracy) will occur. In order to prevent uneven illumination, a document area of about several mm in the sub scan direction need be illuminated in proper light amount.\nParticularly, a light receiving portion (document read area) of a digital copying machine or an image scanner is very narrow in the sub scan direction, about 0.1 mm, for example. Because of this, the document read area need be illuminated in proper amount such that the center of a luminance distribution curve comes at the document read area; otherwise, luminance of the area is substantially decreased.\nThere is a demand for a document illuminating system for use in a digital copying machine or an image scanner to illuminate the document read area with even luminance even when the center of illumination shifts from the document read area, so that a wide area of the image sensor in the sub scan direction receives light in proper amount with even luminance distribution.\nPreferably, the luminance distribution is to include, near a maximal value, an even portion in a width required to read a document with a fluctuation due to a mechanical error or the like added.\nThe even portion refers to a portion of the luminance distribution in which a data change rate of electric signals can fall within a practically allowable range by correcting the level of electric signals photo-electrically converted from an image by use of a signal processing circuit provided after the image sensor. For generating monochrome images, an input data change rate of about 30% is allowable by electrically correcting (amplifying) image signal values (gain adjustment). However, for generating color images, the input data change rate has to be about 12% or less since correction of color unbalance of three ROB colors need be taken into consideration rather than the correctable range of gain adjustment.\nHowever, References 1 to 3 cannot achieve even luminance distribution and need to be improved."}
-{"text": "The recycling of old rubber tires presents a very difficult problem. In the first place, there is considerable difficulty in separating embedded metal parts from the rubber of the tire before the tire runs through the pulverizing machine. Moreover, there is a possibility of reusing the rubber only when the rubber is reduced practically to powder form because it is only in this form that old rubber can be used as an additive in the rubber industry for the production of new rubber ware.\nIn order to comminute old rubber tires, cutting and shredding tools have been used. According to DE 29 11 251 C2, the tire is clamped in a clamping device, with three outwardly slidable pins, whereupon a tangential cut is made with rotating disk knives to remove the tread with the side strips from the metal inlay containing beads and subsequently other circular disk-form knives are used to cut the rubber into strips which are then fed into cross cutting apparatus. The rubber pieces thus obtained required further comminuting.\nOther old rubber tire disintegrating machines are shredders which, with rotating knife cylinders, attack the circumference of the tread and comminute the tread as well as the side strips to chips (SU 36 93 894 A1). A machine according to SU 13 88 294 A1 works in a similar manner. Here the tire is gripped on both sides by saucers and is stressed and pressed so that the tread is bent and deformed and in this narrowed form is subjected to the knives of a shredding cylinder. Comminuting is effected in a similar manner by the comminuting machine of DE 37 04 725 A1, in which the comminuting tools arranged on the rotating comminuting cylinder are small, hard plates, and the comminuting tools themselves are arranged in a particular form.\nThe old rubber pieces obtained from all of these comminuting processes are technically not yet usable in the rubber industry as an additive. The rubber pieces thus obtained are so large, of such a dissimilar shape and so nonhomogeneous that further working and comminuting presents great difficulty.\nRubber scrap is also produced in the production of new rubber ware when, through production error or mixing error or through use of unsuitable mixtures, damaged goods are produced.\nApparatus for further working such rubber scrap is known through DD 265 855 A1. With this known apparatus, an extruder with several work zones is connected by a radial discharge pipe with a mill having grinding disks defining a conical grinding gap. Different work zones are formed in the extruder through different formation of different screw sections. Bringing the compressed rubber scrap from the extruder into the grinding space is difficult and the grinding of the rubber is problematic.\nFor rubber is an extremely difficult to comminute material, which is in many respects nonhomogeneous. Through different exposure to the sun's rays, and different loading during the life of an automobile tire, as well as different life spans, the old rubber scrap to be worked on is differently aged. The pieces of old rubber scrap have strongly aged portions, which are relatively easy to pulverize. Other, less aged parts, on the contrary, still have a high degree of elasticity and can hence be comminuted by grinding only with great difficulty or not at all. In most cases, strongly aged portions are integral with less aged portions. They thereby present particular difficulty in comminuting. This difficulty is increased by the presence of foreign material in the form of metal parts in the rubber scrap which are not removable prior to grinding because they are in part surrounded by the rubber and in part the rubber is vulcanized to them. These metal parts have a destructive effect on the grinding disks."}
-{"text": "The invention relates to a process for the preparation of mono- and oligoisocyanates by reacting primary amines with phosgene in the presence of a catalyst.\nIsocyanates are industrial products which have a large number of uses in the field of polyurethane plastics. However, certain isocyanates are also used in the preparation of pharmaceutical active ingredients.\nThe synthesis of isocyanates by reacting amines with phosgene has been known for some time. In principal, two processes are described in the literature, one of which is carried out at atmospheric pressure and the other is carried out at increased pressure. Phosgenation under increased pressure is disadvantageous since it requires much more complex industrial apparatus to control the increased safety risk, the release of phosgene.\nFor sulfonyl isocyanates, U.S. Pat. No. 3,371,114 and U.S. Pat. No. 3,484,466 disclose a preparation process at atmospheric pressure in which a solution of a sulfonylamide and an isocyanate as catalyst in an inert solvent is reacted with phosgene. In the process, the corresponding sulfonylurea is formed as an intermediate, which reacts with phosgene to give the desired sulfonyl isocyanate.\nAlkyl and aryl isocyanates are usually prepared by the phosgenation process, described, for example, in Houben-Weyl, Methoden der organischen Chemie [Methods in Organic Chemistry], 4th Edition, Volume E4, pages 741-751, Georg Thieme Verlag Stuttgart, 1983, from the corresponding amines in two phases at atmospheric pressure. In the first phase, the cold phosgenation, the amine is reacted with an excess of phosgene in very dilute solution and at low temperatures to give the corresponding carbamyl chloride, from which, in the second phase at elevated temperature, the hot phosgenation, the isocyanate forms. Aliphatic and cycloaliphatic primary amines are more difficult to phosgenate because of their increased basicity compared with aromatic amines, and lead to an increased formation of byproducts. A disadvantage of these processes is, in addition to the fact that the phosgenation is carried out in two phases, the formation of an intermediate solids suspension of sparingly soluble carbamyl chloride and amine hydrochloride, which in turn renders an increased dilution of the reaction medium necessary in order to prevent deposits and blockages of parts of the equipment. Because of the accumulation of solids which occurs, this process cannot be carried out continuously at atmospheric pressure. Furthermore, symmetrically N,Nxe2x80x2-substituted urea forms as a byproduct, the formation of which can only be suppressed at the expense of drastically reduced space-time yields.\nAliphatic and cycloaliphatic amines are frequently used in the form of their salts in the cold/hot phosgenation. However, these salts are sparingly soluble in the reaction medium, meaning that additional reaction stages and very long reaction times are necessary.\nFurthermore, it is known from GB 1 114 085, U.S. Pat. No. 3,492,331 and H. Ulrich, Chemistry and Technology of Isocyanates, Wiley and Sons, 1996, pages 328-330, to optimize the reaction of primary amines with phosgene by the addition of catalysts such as dimethylformamide, phenyltetramethylguanidine, 2,4,6-trimethylpyridine or carbodiimidazole. Some of these catalysts must be used in equimolar amounts and form sparingly soluble salts under the reaction conditions.\nIt is an object of the present invention to provide a process, which can be used for aliphatic, cycloaliphatic, araliphatic and aromatic primary amines, for the preparation of the corresponding mono- and oligoisocyanates using phosgene which can be carried out either continuously or batchwise at atmospheric pressure, does not have the above disadvantages and provides the corresponding mono- and oligoisocyanates in good yields and high selectivities.\nWe have found that this object is achieved by a process for the preparation of aliphatic, cycloaliphatic, araliphatic and aromatic mono- and oligoisocyanates by phosgenation of the corresponding primary amines at atmospheric pressure, in which\na) a catalytic amount of a monoisocyanate (isocyanate a) is introduced into an inert solvent together with phosgene,\nb) the primary amine is added, and\nc) the resulting reaction mixture is reacted with phosgene.\nThe net equation underlying the process is given in scheme 1 below. \nThe process according to the invention has the advantage that it can be used for a large number of amines. The phosgenation is carried out according to the invention, while avoiding a division into cold and hot phosgenation, in a narrow temperature interval and at atmospheric pressure, the intermediate formation of sparingly soluble suspensions being avoided. The desired isocyanate is formed in the process, with complete conversion of the amine, in high yields and high selectivity in significantly shortened reaction times without symmetrically substituted N,Nxe2x80x2-urea being formed from the amine as a byproduct. Since the formation of urea from the amine is not observed, it is possible by means of the process according to the invention to significantly increase the concentration of the amine in the reaction solution and thus the space-time yields. In addition, it is advantageous that the process according to the invention can be carried out either batchwise or continuously since there is no accumulation of solids.\nThe process according to the invention gives aliphatic, cycloaliphatic, araliphatic and aromatic mono- and oligoisocyanates of the formula I\nR1xe2x88x92N=C=Oxe2x80x83xe2x80x83(I)\nThe radical R1 in formula I corresponds to the radical R1 in formula IV of the amines used in the process according to the invention, which are discussed later and to which reference is made here. Preference is given to preparing mono- and diisocyanates by the process of the invention. Of lesser importance in practice, but preparable in principle are isocyanates having 3 and more isocyanate groups.\nThe sole catalyst used is a monoisocyanate of the formula II (isocyanate a)\nR2xe2x88x92N=C=Oxe2x80x83xe2x80x83(II)\nor mixtures thereof, in which R2 is aliphatic, cycloaliphatic, aromatic or araliphatic radicals. These can be substituted by heteroatoms, or their carbon chains can be interrupted by heteroatoms, such as oxygen and sulfur. The radical R2 must, however, be inert toward phosgene, thus excluding radicals which carry NH, OH and SH groups. The aliphatic radicals can be arbitrarily branched or unbranched, saturated or unsaturated. They contain 3 to 30 carbon atoms, preferably 3 to 10 carbon atoms. Examples of aliphatic radicals are methyl, ethyl, propyl, n-butyl, isobutyl and sec-butyl.\nSuitable cycloaliphatic radicals are those which have 3 to 20 carbon atoms, preferably 3 to 10 carbon atoms, such as, for example, cyclopentyl and cyclohexyl.\nThe aromatic radicals can be unsubstituted or arbitrarily substituted by alkyl or aryl substituents or heteroatoms. Preference is given to the aromatic radicals which are mono- or disubstituted. Examples of aromatic radicals are phenyl, chlorophenyl, o-, m- and p-tolyl.\nSuitable araliphatic radicals are radicals having 7 to 12 carbon atoms, such as, for example, benzyl, although preference is given to a radical of the formula III\n\nin which R3 and R4 can be identical or different and can be hydrogen or an aliphatic, cycloaliphatic, araliphatic or aromatic radical, in which case these may be substituted by heteroatoms, or their carbon chains can be interrupted by heteroatoms, such as oxygen and sulfur. The radicals R3 and R4 must, however, be inert toward phosgene, thus excluding radicals carrying NH, OH and SH groups. The aliphatic radicals can be arbitrarily branched or unbranched, saturated or unsaturated. They contain 1 to 30 carbon atoms, preferably 1 to 10 carbon atoms. Examples of aliphatic radicals are methyl, ethyl, propyl, n-butyl, isobutyl and sec-butyl.\nSuitable cycloaliphatic radicals are those which have 3 to 20 carbon atoms, preferably 3 to 10 carbon atoms, such as, for example, cyclopentyl and cyclohexyl.\nThe aromatic radicals can be unsubstituted or arbitrarily substituted by alkyl or aryl substituents or heteroatoms. Preference is given to the aromatic radicals which are mono- or disubstituted. Examples of aromatic radicals are phenyl, chlorophenyl, o-, m- and p-tolyl. R5 is fluorine, chlorine, bromine or a C1- to C4-alkyl chain which may optionally be interrupted by a heteroatom, for example oxygen or sulfur. Preference is given to the monoisocyanate (isocyanate a), a linear aliphatic alkyl isocyanate, particularly preferably a linear aliphatic alkyl isocyanate having 3 to 10 carbon atoms, where the alkyl chain may optionally be branched, particularly preferably n-butyl isocyanate. In the case of the preparation of monoisocyanates by the process according to the invention, it is also possible to use the monoisocyanate desired as product in catalytic amounts as isocyanate a. It is, however, preferable to use an isocyanate a which is not identical to the product (isocyanate of the formula (I)) as catalyst.\nThe isocyanate a is generally used in a catalytic amount of from 0.01 to 25 mol %, preferably 0.5 to 20 mol %, particularly preferably 1 to 15 mol %, based on the amine.\nIn step a) of the process according to the invention, the isocyanate a is introduced into a solvent which is inert toward phosgene. Preferred solvents are hydrocarbons. Particular preference is given to mono- or polysubstituted aromatic hydrocarbons, such as toluene, o-, m- or p-xylene, ethylbenzene, chlorobenzene or 1,2-, 1,3- or 1,4-dichlorobenzene or mixtures thereof. Particular preference is given to xylenes, chlorobenzene and dichlorobenzenes as solvent.\nThe process according to the invention is carried out at atmospheric pressure and generally at a reaction temperature, which is largely constant over the three steps of the process a), b) and c), of from 20 to 200xc2x0 C., preferably from 20 to 150xc2x0 C., particularly preferably from 50 to 120xc2x0 C.\nThe monoisocyanate (isocyanate a) dissolved in the inert solvent is, in step a) of the process according to the invention, heated to the reaction temperature and mixed with an excess of phosgene, based on the monoisocyanate. The molar ratio between phosgene and isocyanate a is preferably 100:1 to 5:1. Particular preference is given to using phosgene in an excess of 20:1 to 5:1, based on isocyanate a.\nThe phosgene can be used in the process according to the invention in steps a) and c) in gaseous form or in condensed form, preference being given to introducing the phosgene in gaseous form.\nIn the process according to the invention, suitable primary amines are amines of the formula IV\nR1xe2x88x92NH2xe2x80x83xe2x80x83(IV)\nwhich are obtainable by methods known from the literature or as a commercially available product.\nR1 stands for aliphatic, cycloaliphatic, aromatic or araliphatic radicals. These can be substituted by heteroatoms, or their carbon chains can be interrupted by heteroatoms, such as oxygen and sulfur. Radicals which contain OH and SH groups are excluded. Substitution by amino groups, on the other hand, is possible and leads to di- or oligoisocyanates, preference being given to diisocyanates. The aliphatic radicals can be arbitrarily branched or unbranched, saturated or unsaturated. They contain 3 to 30 carbon atoms, preferably 3 to 10 carbon atoms. Examples of aliphatic radicals are propyl, butyl, pentyl and hexyl. Suitable cycloaliphatic radicals are those having 3 to 20 carbon atoms, preferably those having 3 to 10 carbon atoms, such as cyclopentyl and cyclohexyl.\nThe aromatic radicals can have 4 to 10 carbon atoms, can be heteroaromatic radicals, and unsubstituted or arbitrarily substituted by alkyl or aryl substituents or heteroatoms. The aromatic radicals are preferably mono- or disubstituted. Examples of aromatic radicals are phenyl and chlorophenyl, o-, m-, p-tolyl, furfuryl, thiophenyl and pyrrolyl.\nSuitable araliphatic radicals are radicals having 7 to 12 carbon atoms, such as, for example, benzyl, although preference is given to a radical of the formula V \nin which R6 and R7 can be identical or different and can be hydrogen or an aliphatic, cycloaliphatic or aromatic radical, in which case these can be substituted by heteroatoms, or their carbon chains can be interrupted by heteroatoms, such as oxygen and sulfur. The radicals R6 and R7 must, however, be inert toward phosgene, meaning that radicals carrying NH, OH and SH groups are excluded. The aliphatic radicals can be arbitrarily branched or unbranched, saturated or unsaturated. They contain 1 to 30 carbon atoms, preferably 1 to 10 carbon atoms. Examples of aliphatic radicals are methyl, ethyl, propyl, n-butyl, isobutyl and sec-butyl.\nSuitable cycloaliphatic radicals are those which have 3 to 20 carbon atoms, preferably 3 to 10 carbon atoms such as, for example, cyclopentyl and cyclohexyl.\nThe aromatic radicals can be unsubstituted or arbitrarily substituted by alkyl or aryl substituents or heteroatoms. Preference is given to the aromatic radicals which are mono- or disubstituted. Examples of aromatic radicals are phenyl, chlorophenyl, o-, m- and p-tolyl. R8 is a halogen atom (F, Cl, Br) or a C1- to C4-alkyl chain, which may optionally be interrupted by a heteroatom, for example oxygen or sulfur. Primary amines of the formula IV which can be used are also amines in enantiomerically pure, optically active form or mixtures of the enantiomers. The process according to the invention advantageously leads during the reaction and the work-up to only a low degree of racemization, which is generally below 2%.\nIn the process according to the invention, preference is given to using cycloaliphatic, araliphatic and aromatic amines, in particular those which are further processed in industrial processes, such as, for example, hexamethylenediamines, cyclohexylamine, isophoronediamine, toluylenediamine, aniline, and optically active amines, such as, for example, R-(+)- and S-(xe2x88x92)-phenylethylamine.\nIn step b) of the process according to the invention, the primary amine is added to the solution of the monoisocyanate (isocyanate a), saturated with phosgene, preferably in an inert solvent. This can be carried out, for example, in the presence of a stream of inert gas, for example nitrogen. Preferred solvents are the solvents used in step a) of the process according to the invention. Particular preference is given to using the same solvent as in step a). The resulting reaction mixture is then reacted with phosgene in step c) of the process according to the invention.\nThe molar ratio between primary amine and the total amount of the amount of phosgene introduced in steps a) and c) is 1:10 to 1:1, particular preference being given to introducing phosgene in an amount of from 120 to 180 mol %, based on the amine.\nThe process according to the invention can in principle be carried out both continuously and batchwise, preference being given to a continuous operation. The process can be carried out in any desired apparatus which is suitable for a reaction with phosgene. For example, a phosgenation apparatus is used which comprises an attached carbonic acid condenser or a downstream battery of high-efficiency condensers for the condensation of phosgene.\nFor the preferred continuous procedure, a phosgenation apparatus is used which comprises a main reactor and after-reactor, a stripping column and, connected downstream therefrom, a condenser. In the case of the continuous procedure, a mixture of the desired isocyanate of the formula I, the catalytic amounts of the isocyanate of the formula II (isocyanate a) and phosgene in the inert solvent is preferably initially introduced.\nAfter an after-reaction, excess phosgene and solvent are generally stripped out at temperatures from 30 to 80xc2x0 C., preferably 40xc2x0 C. to 60xc2x0 C., condensed in the condenser and returned together to the after-reactor.\nThe reaction product is worked up using methods known to the person skilled in the art. Preferably, the desired isocyanate of the formula I is isolated by fractional distillation. The isocyanate a used as catalyst is preferably recovered by distillation and can be reused for further batches. The desired mono- and oligoisocyanates can be obtained in the process in yields of generally 60 to 95%, based on the primary amine used.\nThe examples below illustrate the invention further."}
-{"text": "The present invention relates generally to an open superconductive magnet used to generate a uniform magnetic field as part of a magnetic resonance imaging (MRI) diagnostic system, and more particularly to such a magnet having a support structure for a high magnetic field strength.\nMRI systems employing superconductive or other type magnets are used in various fields such as medical diagnostics. Open magnets employ two spaced-apart superconductive coil assemblies (or other magnet pole assemblies) with the space between the assemblies allowing for access by medical personnel for surgery or other medical procedures during MRI imaging. The patient may be positioned in that space or also in the bore of annular coil assemblies. The open space helps the patient overcome any feelings of claustrophobia that may be experienced in a closed magnet design. Known open magnet support systems include \"C\"-shaped supports with the superconductive coil assemblies (or other magnet pole assemblies) being attached to the open top and bottom ends of the \"C\". Other known support systems include structural posts whose ends are attached to the two superconductive coil assemblies (or other magnet pole assemblies) including attachment to the coil forms (which hold the superconductive coils) of such superconductive coil assemblies. These support systems may be adequate for low field strength magnets, but a more robust support system is desired when the open magnets have a high field strength, including one Tesla or higher, where the attraction forces between the toroidal-shaped superconductive coil assemblies can be in excess of 160,000 pounds."}
-{"text": "This application is a National Stage of International Application No. PCT/EP2007/005557, filed Jun. 23, 2007, which claims priority under 35 U.S.C. \u00a7119 to German Patent Application No. 10 2006 038 124.6, filed Aug. 14, 2006, the entire disclosure of which is herein expressly incorporated by reference.\nThe present invention relates to a restraint system for a motor vehicle.\nIt is well known for motor vehicles to protect occupants from injuries by inflating an air bag in a short period of time, to catch the occupant who is displaced in a forward direction. For this purpose, it is important that the air bag quickly reaches its effective volume, and it is therefore necessary to produce relative large amounts of gas (e.g., in a pyrotechnical manner) and to introduce them into the gas bag. Recently, for dimensioning the gas volume to be added, occupant/and or vehicle parameters have been evaluated, so as not to unfold the gas bag to its full size for example in the case of a so-called \u201cout of position\u201d of the occupant. However, the necessity for adjustment to all possible load conditions and the relatively cost-intensive control and regulation technique have made this feature problematic.\nIntensified efforts have been made, therefore to develop so-called self-regulating systems which can adjust automatically to the corresponding load conditions. For example, venting apertures have been provided, which close automatically when the inner pressure of the gas bag becomes too large.\nGerman patent document DE 2 302 737 discloses a restraint system comprising a two-layer gas bag, where the gas is guided only between the two layers, so that a complete gas cushion does not result; rather, a spherical annular supporting structure. The unfolding of the support structure thereby occurs in the transverse direction of the abutting gas cushions. On the other hand, European patent document EP 0589 059 B1 shows furthermore, that it is necessary during the unfolding of the two-layer gas bag, to suck ambient air into the interior so as to overcome the negative pressure.\nBoth of these systems have in common the feature that less gas volume is necessary to unfold the gas bag to its full size, due to the gas bag's being formed with two layers. The temperature and the pollution can thereby be reduced.\nFinally, UK patent document GB 1 420 226 A discloses a restraint system for a motor vehicle, where a tubular supporting structure is provided in the interior of the two-layer gas cushion, the longitudinal extension of which exceeds the dimensions of the transverse extension in the inflated (that is, active) state. The supporting structure unfolds due to its particular geometry mainly in the direction of its longitudinal extension.\nAn automatic adjustment of the gas bag dependent on the respective load condition (e.g., dependent on the respective occupant or his or her position) is not described.\nIt is therefore an object of the present invention to improve the generic restraint system in such a manner that its size adjusts to the respective load condition.\nThis and other objects and advantages are achieved by the restraint system according to the invention, which when it impacts upon an obstacle during unfolding, does not achieve the stability it would have when fully unfolded and/or the final volume or the final extension, due to the particular geometric form of the supporting structure (that is, due to its longitudinal extension exceeding its cross section). The supporting structure has a different stability during unfolding, due to its longitudinal formation. While the supporting structure is initially rather unstable (that is, can be easily hindered during unfolding), the completely unfolded supporting structure achieves full stability (that is restraining force). This means that, if the supporting structure impacts upon an obstacle during the unfolding, as is the case for example with occupants leaning forward (out of position), unfolding can be stopped more easily or deflected due to the supporting structure which is still unstable. These small forces effect a low pressurization of the occupant.\nAs used herein, the term \u201csupporting structure\u201d refers to a structure which is similar to a skeleton (in contrast to a conventional spherical gas bubble\u2014a gas bag), and which comprises a restraint effect in the fully unfolded state comparable to the conventional gas bag. The provision of a substantially more complex inflatable supporting structure according to the invention, for example a branched tree structure not only reduces the necessary gas volume, but also diminishes the force peaks acting on the occupant during unfolding, if he is effectively \u201cin the way\u201d of the unfolding supporting structure. In contrast, the gas amount in the conventional gas bubble is on a substantially higher level from the start of the activation until the full unfolding, so that obstacles in the unfolding path are put under more pressure independent of the unfolding state.\nThe supporting structure can preferably be filled with gas, and it is unimportant whether it is filled with gas generated in a pyrotechnical manner or with gas from a pressure vessel or the ambient air. The ambient air can be guided into the carrier volume enclosed by the supporting structure by suction through the negative pressure occurring during the unfolding. The gas volume can be reduced further in this manner.\nIf several supporting structures are provided, which are connected to one another or to the environment in a fluidic manner, a restraint system can be formed, which is constructed in a net-like, tree-like or supporting frame-like manner. Thus, supporting structures proceeding in the longitudinal direction of the vehicle can be connected to supporting structures proceeding transversely to the transverse direction of the vehicle. The unfolding of the transverse supporting structures thereby nevertheless takes place in the direction of the longitudinal extension of the respective supporting structure.\nThe supporting structures can be connected by flexible sheets, in particular flexible textile sheets. In this case, the flexible sheet can be fixed between two adjacent supporting structures.\nIf several supporting structures are connected to one another by a flexible sheet, the gas can be used for restraining action in the resulting enclosed carrier volume enclosed. The gas in the carrier volume serving to restrain occupants can be heated or supplemented by means of a heating device (e.g., an ignition tablet or a small gas generator step, which develops heat with a possibly very low gas volume). As a result, the gas volume or the inner pressure, and thus the restraint action, increases therewith correspondingly.\nIf a supporting structure impacts upon an obstacle during unfolding, the gas flowing into the supporting structure can be distributed to other adjacent supporting structures or into the environment, by gas redistribution components. The inner pressure generated by the flow in the supporting structure can also be reduced by increasing its cross section. This can take place in such a manner that tear seams break down at a certain pressure, so that the supporting structure can increase in its transverse direction. A pressure-relief valve can also be provided, which opens when an interior pressure is reached which is too high.\nWith conventional compact gas bags, it was previously necessary to refrain from directing venting apertures towards the occupant, because gas temperatures are reached which are too high. However, with the system according to the invention, as cold ambient air flows within the carrier volume, and not a pyrotechnically produced gas, and venting apertures can be directed towards the occupant by the reduced temperature. The air bag dampening can thus be adapted for various environmental conditions, by sealing a venting aperture cross section between occupant and air bag.\nThus, this contact surface and its sealing with persons having a higher volume (and usually a higher weight) is larger, so that a stronger restraint action is achieved hereby. With more severe accidents, the contact surface and thus the restraint action is also increased by a stronger immersion of the occupant into the air bag. This principle also permits a variable air bag dampening for belted and unbelted occupants, because in the unbelted state, the occupant will be immersed into the air bag with less force.\nIf the venting apertures are formed as perforations (that is, many small apertures), as for example with a textile net, a projection area corresponding to the measurements of the occupant can be closed. The reproducibility of the results increases on average with many small venting apertures.\nIn a preferred embodiment, at least one of the supporting structures extends at least partially in an overlap of a hard structure, as for example a supporting column. It is also possible to provide a structure which unfolds between the occupants as an interaction bag.\nSo as to be able to manufacture a complex supporting structure with its inflatable components and the flexible sheets, the one-piece woven technique is recommended. This technique distinguishes itself in that, on one and the same device, a double layer can be manufactured for the inflatable structures, and one layer for the flexible sheets, or a three-dimensional structure can be woven.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
-{"text": "Polysilicon formed on a field oxide layer may be used as a resistor element for use in the semiconductor integrated circuit. In such a case, a plurality of polysilicon members having the same shape and the same characteristics may be formed as a resistor element unit for the purpose of providing a desired resistance value. The number of polysilicon members connected in series or in parallel is then adjusted to adjust the resistance value. The size and shape of such polysilicon members are required to be highly precise. The polysilicon members may thus need to be manufactured under the same process conditions. Namely, the conditions of photolithography and etching performed during the wafer process may need to be maintained constant with respect to each polysilicon member.\nWhen a plurality of polysilicon members having the same shape are arranged in matrix form at constant intervals, polysilicon members situated in the interior of the matrix are formed under the same process conditions because each of these polysilicon members has identical positional relationships with the surrounding adjacent polysilicon members. The polysilicon members situated at the periphery of the matrix, however, have adjacent polysilicon members only on one side thereof, and, thus, are subjected to different process conditions than the conditions applied to the interior polysilicon members which are surrounded by adjacent polysilicon members on the four sides thereof. Accordingly, the polysilicon members situated at the periphery are provided, as dummy members unused for actual circuit operations, only for the purpose of achieving constant process conditions for the interior polysilicon members.\nFIG. 1 is a plan view illustrating an example of a related-art resistor-element unit. FIG. 2 is a cross sectional view of the resistor-element unit illustrated in FIG. 1 as taken along a line A-A\u2032.\nAs illustrated in FIG. 1 and FIG. 2, a resistor-element unit 10 includes a P-type substrate 13, a field oxide layer 14 formed on the P-type substrate 13, and a plurality of polysilicon members 11 and 12 formed on the field oxide layer 14. The polysilicon members 11 and 12 are arranged in matrix form at constant intervals which are an interval \u201ca\u201d in the vertical direction and an interval \u201cb\u201d in the horizontal direction. The polysilicon members 12 are dummy polysilicon members situated at the periphery of the matrix, and the polysilicon members 11 are resistor-element polysilicon members situated in the interior of the matrix.\nThe resistor-element polysilicon members 11 situated in the interior of the matrix are formed under the same process conditions because each of these polysilicon members has identical positional relationships with the surrounding adjacent polysilicon members. Accordingly, the resistor-element polysilicon members 11 have the same resistor-element characteristics. A desired resistor value can thus be obtained by connecting a desired number of resistor-element polysilicon members 11 in series or in parallel.\nThe dummy polysilicon members 12 situated at the periphery of the matrix, however, have adjacent polysilicon members only on one side thereof, and, thus, are subjected to different process conditions than the conditions applied to the interior resistor-element polysilicon members 11 which are surrounded by adjacent polysilicon members on the four sides thereof. Accordingly, the dummy polysilicon members 12 situated at the periphery are formed, as dummy members unused for actual circuit operations, only for the purpose of achieving constant process conditions for the interior polysilicon members.\nIn the resistor-element unit 10 formed as described above, the dummy polysilicon members 12 situated at the periphery are not used as part of an actual operating circuit, which means that there are needless elements and areas provided in the semiconductor integrated circuit chip. This gives rise to problems such as an increase in chip size and a cost increase."}
-{"text": "German Patent Application No. 28 29 057, discusses a fuel-injection system for supplying fuel to a mixture-compressing internal combustion engine having external ignition as a function of operating parameters. The fuel-injection system includes a metal fuel-distributor line, which, via at least one branch line, is connected to at least one fuel injector, the branch line being embodied as a metal tube and connected to the fuel injector by way of a threaded connection. Easily bendable metal is used as material for the branch line. Situated between the threaded connection on the branch line and the fuel injector are thin-walled metal bellows, which compensate for a lateral offset between the beginning point of the branch line on the fuel-distributor line and the fitting position of the fuel injector, and which damp the operating noises originating in the fuel injector.\nDisadvantageous in the fuel-injection system from German Patent Application No. 28 29 057 is that the screw connection is not secured against an automatic loosening. Due to the vibrations of the internal combustion engine during operation, a screwed connection rigidly connected to the internal combustion engine is at a higher risk of coming unscrewed.\nHowever, especially in the case of directly injecting fuel injectors and at the high pressures required in this context, the screwed connection is safety-relevant and should not come unscrewed under any circumstances. The known fuel-injection system provides no suggestion for securing a screwed connection."}
-{"text": "There is known a technology for detecting the occurrence of a misfire in each of cylinders of an engine. See, Japanese Laid-Open Patent Publication No. 2009-280082 (hereinafter referred to as \u201cJP2009-280082A\u201d). According to JP2009-280082A, when a management ECU (117) judges that a misfire has occurred in an internal combustion engine (107), the management ECU turns on a warning lamp (125) (see paragraph [0027]).\nThere also is known a technology for detecting an abnormal compression pressure in each of cylinders of an engine. See, Japanese Laid-Open Patent Publication No. 2004-019465 (hereinafter referred to as \u201cJP2004-019465A\u201d). According to JP2004-019465A, while a fuel system and an ignition system are inactivated, the engine is cranked in order to rotate a crankshaft (1a), and a rotational variation of the engine is detected using a difference between instantaneous rotational speeds at preset crankshaft angles in compression strokes of the cylinders. An abnormal compression pressure is detected based on the rotational variation (see claims 1 through 6)."}
-{"text": "Carburetors are devices that can be used to mix fuel and air to power combustion engines typically including gasoline powered internal combustion spark ignited engines. A carburetor may include a fuel metering system that helps to control the amount of fuel supplied to air flowing through a mixing passage or main bore of the carburetor for mixing the fuel with air and supplying the mixture to the engine. Some metering systems employ a diaphragm that oscillates or reciprocates during operation to open and close a metering valve admitting fuel to a chamber from which it is supplied to the passage for mixing with air. In use, the large number of cycles experienced by such a diaphragm which typically physically interacts with other components of the metering system such as a valve actuating lever, and continuous exposure to solvent containing fuels, can result in a harsh operating environment that causes wear, degradation and ultimately failure of the diaphragm. In a gasoline powered spark ignited internal combustion so-called small engine the diaphragm must fully open the valve when subjected to only a small pressure differential which is typically a maximum negative pressure of \u22120.9956 kPa or \u22120.1444 pounds per square inch and usually about \u22120.50 kPa or \u22120.0725 pounds per square inch (psi). This very small differential operating pressure also requires that the portion of the diaphragm within a fuel metering chamber be very flexible particularly since such diaphragm may have a surface area within the metering chamber in the range of about 0.5 square inch to 1.0 square inch."}
-{"text": "1. Field of the Invention\nThis invention relates generally to the field of memory interface design and, more particularly, to cache design in a graphics system.\n2. Description of the Related Art\nA computer system typically relies upon its graphics system for producing visual output on the computer screen or display device. Early graphics systems were only responsible for taking what the processor produced as output and displaying that output on the screen. In essence, they acted as simple translators or interfaces. Modern graphics systems, however, incorporate graphics processors with a great deal of processing power. They now act more like coprocessors rather than simple translators. This change is due to the recent increase in both the complexity and amount of data being sent to the display device. For example, modern computer displays have many more pixels, greater color depth, and are able to display images that are more complex with higher refresh rates than earlier models. Similarly, the images displayed are now more complex and may involve advanced techniques such as anti-aliasing and texture mapping.\nSince graphics systems typically perform only a limited set of functions, they may be customized and therefore far more efficient at graphics operations than the computer's general-purpose central processor. While early graphics systems were limited to performing two-dimensional (2D) graphics, their functionality has increased to support three-dimensional (3D) wire-frame graphics, 3D solids, and now includes support for three-dimensional (3D) graphics with textures and special effects such as advanced shading, fogging, alpha-blending, and specular highlighting.\nWith each new generation of graphics system, there is more image data to process and less time in which to process it. This consistent increase in data and data rates places additional burden on the memory systems that form an integral part of the graphics system. One example of a memory sub-system defining the upper limit of overall system performance may be the texture buffer of a graphics system. Certain graphics applications such as 3D modeling, virtual reality viewers, and video games may call for the application of an image to a geometric primitive in lieu of a procedurally generated pattern, gradient or solid color. In these applications, geometric primitives carry additional mapping data (e.g., a UV, or UVQ map), which describes how the non-procedural data is to be applied to the primitive. To implement this type of function, a graphics system may employ a texture buffer to store two dimensional image data representative of texture patterns, \u201cenvironment\u201d maps, \u201cbump\u201d maps, and other types of non-procedural data.\nDuring the rendering process, the mapping data associated with a primitive may be used to interpolate texture map addresses for each pixel in the primitive. The texture map addresses may then be used to retrieve the portion of non-procedural image data in the texture buffer to be applied to the primitive. In some cases (e.g., photo-realistic rendering) a fetch from the texture buffer may result in a neighborhood or tile of texture pixels or texels to be retrieved from the texture buffer and spatially filtered to produce a single texel. In these cases, four or more texels may be retrieved for each displayed pixel, placing a high level of demand on the texture buffer. Thus, poor performance of the texture buffer is capable of affecting a cascading degradation through the graphics system, stalling the render pipeline, and increasing the render or refresh times of displayed images.\nIn other words, accesses of graphics data, such as texture map data, must be performed very quickly. Therefore, one goal of a graphics system is to improve the speed and efficiency of memory accesses of texture maps from a texture memory. One common method is to use a texture memory cache to improve the speed of accesses of texture maps from the texture memory. The design of texture memory systems, including texture cache memory systems, plays a significant role in the implementation of new generation graphics systems.\nTexture mapping hardware often needs to process multiple pixels in the same cycle in order to maximize performance. While these pixels are independent, they typically exhibit some degree of spatial coherence. A high performance texture cache would take advantage of this coherence without imposing ordering or synchronization restrictions between pixels. Texture mapping is generally a read-only operation. Consequently, latency is typically not a factor when considering proper operations, but bandwidth does affect performance. In contrast, the performance of microprocessor instruction and data caches is typically affected by latency, not bandwidth.\nIn a texture mapping graphics system that processes multiple pixels, typically in multiple parallel pipelines, it is common for neighboring pixels to request overlapping texture data. Traditionally, cache memory in such a system requires one port per texel per pipeline, which is generally expensive to implement, or may suffer performance loss from having too few read-ports. Alternately, multiple cache memories may be built into the texture mapping graphics system. This solution is typically expensive to implement as well.\nTherefore, new texture cache memory systems and methods are desired to optimize texture caching by reducing the number of cache read-ports necessary to support the simultaneous reading of texels for multiple pixels, taking advantage of the spatial locality of the texel accesses. More generally, improved cache memory systems are desired in various different applications, including graphics applications."}
-{"text": "1. Field\nThe present invention generally relates to induction RF fluorescent light bulbs, and more specifically to reduction of electromagnetic interference from an induction RF fluorescent light bulb with a ferromagnetic core.\n2. Description of Related Art\nDischarge lamps create light by exciting an electrical discharge in a gas and using that discharge to create visible light in various ways. In the case of fluorescent lamps the gas is typically a mixture of argon, krypton and/or neon, plus a small amount of mercury. Other types of discharge lamps may use other gasses. The gas is contained in a partially evacuated envelope, typically transparent or translucent, typically called a bulb or arc tube depending upon the type of lamp.\nIn conventional discharge lamps electrically conductive electrodes mounted inside the bulb or arc tube along with the gas provide the electric field used to drive the discharge.\nUse of electrodes can create certain problems. First, the discharge is typically designed to have a relatively high voltage in order to minimize losses at the electrodes. In the case of fluorescent lamps, this may lead to long, thin lamp structures, which function well for lighting office ceilings, but are not always a good fit for replacing conventional incandescent lamps. Fluorescent lamps designed to replace incandescent lamps, known as compact fluorescent lamps, or CFLs, are typically constructed by bending the long, thin tube, such as into multiple parallel tubes or into a spiral, which is now the most common form of CFLs. A plastic cover shaped like a conventional incandescent lamp is sometimes placed over the bent tubes to provide a more attractive shape, but these covers absorb light, making the lamp less efficient. Bent and spiral tube lamps also have wasted space between the tubes, making them larger than necessary. The use of a cover increases the size further.\nThe use of electrodes can create problems other than shape and size. Electrodes can wear out if the lamp is turned on and off many times, as is typical in a residential bathroom and many other applications. The life of the electrodes can also be reduced if the lamp is dimmed, because the electrodes are preferably operated in a specific temperature range and operation at different power levels can cause operation outside the preferred ranges, such as when operating at lower power, which can allow the electrodes to cool below the specified temperature range.\nIn addition, the long thin shape selected, because it is adapted to allow use of electrodes, tends to require time for mercury vapor to diffuse from one part of the tube to another, leading to the long warm-up times typically associated with many compact fluorescent lamps.\nFinally, the electrodes are normally designed to be chemically compatible with the gas used in the lamp. While this is not usually a concern with typical fluorescent lamps, it can be a problem with other types of discharge lamps.\nOne way to avoid the problems caused by electrodes is to make a lamp that does not use electrodes, a so-called electrodeless lamp. In an electrodeless lamp, the discharge may be driven by, for example, 1) an electric field created by electrodes mounted outside the bulb or arc tube; 2) an electric field created by a very high frequency electromagnetic field, usually in combination with a resonant cavity, or 3) an electric field created by a high frequency magnetic field without the use of a resonant cavity. This latter lamp is called an induction-coupled electrodeless lamp, or just \u201cinduction lamp.\u201d\nIn an induction lamp, a high frequency magnetic field is typically used to create the electric field in the lamp, eliminating the need for electrodes. This electric field then powers the discharge.\nSince induction lamps do not require use of electrodes, they do not need to be built into long thin tubes. In fact, a ball-shaped bulb, such as the bulb used for conventional incandescent lamps, is a preferred shape for an induction lamp. In addition, since induction lamps do not use electrodes, they can be turned on and off frequently without substantial adverse impact on loss of life. The absence of electrodes also means that induction lamps can be dimmed without reducing lamp life. Finally, the ball-shaped lamp envelope allows rapid diffusion of mercury vapor from one part of the lamp to another. This means that the warm-up time of induction lamps is typically much faster than the warm-up time of most conventional compact fluorescent lamps.\nInduction lamps fall into two general categories, those that use a \u201cclosed\u201d magnetic core, usually in the shape of a torus, and those that use an \u201copen\u201d magnetic core, usually in the shape of a rod. Air core induction lamps fall into this latter category. Closed core lamps are usually operated at frequencies generally above 50 kHz, while open core lamps usually require operating frequencies of 1 MHz and above for efficient operation. The lower operating frequency of closed core induction lamps makes them attractive; however, the bulb design required to accommodate the closed core makes them generally unsuitable for replacing standard in incandescent lamps. Open core induction lamps, while requiring higher operating frequencies, allow the design of lamps that have the same shape and size as common household incandescent lamps. This disclosure is primarily is addressed to the open core category of induction lamps.\nIn spite of their obvious advantages, there are very few open core induction lamps on the market today. One reason for the lack of commercially successful products is the cost of the high frequency ballast. Conventional fluorescent lamps, including CFLs, can be operated at frequencies from 25 kHz to 100 kHz, a frequency range where low cost ballast technology was developed in the 1990s, and closed core induction lamps can be operated at frequencies from 50 kHz to 250 kHz, for which the ballasts are only slightly more expensive. However, open core induction lamps typically require operating frequencies of 1 MHz or higher. The United States Federal Communications Commission (FCC) has established a \u201clamp band\u201d between 2.51 MHz and 3.0 MHz that has relaxed limits on the emission of radio frequency energy that may interfere with radio communication services. Cost effective open core induction lamps may preferably have an operating frequency of at least 2.51 MHz.\nThe lack of commercially successful open core induction lamps may be traced to the absence of a low cost ballast that can operate in the 2.51 MHz to 3.0 MHz band while meeting all the requirements of the FCC; that is small enough to fit into a lamp; that has a ballast housing that is the same size and shape as a conventional incandescent lamp; and that can be dimmed on conventional TRIAC dimmers found in homes in certain major markets, such as the U.S. The present disclosure addresses one or more of these issues. Therefore a need exists for improved induction lamps, especially in residential applications."}
-{"text": "Venous, arterial, and body fluid catheters are commonly used by physicians. For example, such catheters may be used to gain access to the vascular system for dialysis, for introducing pharmaceutical agents, for nutrition or fluids, for hemodynamic monitoring, and for blood draws. Alternatively, catheters can be used for drainage of fluid collections and to treat infection. Following introduction into the patient, the catheter is secured to the patient. In conventional practice, the catheter is commonly secured to the patient using an adhesive tape patch or by suturing an attached hub to the patient's skin."}
-{"text": "In recent years, a holey fiber becomes popular as an optical fiber suitable for a long transmission distance. The holey fiber is an optical fiber in which a refractive index of a clad is reduced by holes. It is known that the holey fiber obtains an optical characteristic that cannot be obtained by a conventional optical fiber in which a refractive index of a clad is reduced by an impurity.\nAs described in Patent Literatures 1 through 4, the holey fiber is manufactured by the steps of (1) preparing a columnar base material (hereinafter, referred to as a \u201cpreform\u201d) made from silica glass, (2) forming, in the preform, through holes which are to be holes, and (3) drawing the preform in which the through holes have been formed.\nThe optical characteristic of the holey fiber is influenced by positions of the holes. Accordingly, it is important to form the through holes in predetermined appropriate positions in the preform in order to manufacture an optical fiber having a desired optical characteristic.\nIn the step (3), through holes extending in a direction vertical to end surfaces of the preform are formed by a drilling process. If a machine tool for forming the through holes has low machining accuracy, perforating positions of the through holes are gradually shifted as the perforating is proceeded even if the perforating is started from appropriate positions of one end surface. Accordingly, positions of the through holes are largely shifted at the other end surface, or a through hole is connected to another through hole in the middle of the preform. For this reason, after the through holes are formed, it is necessary to whether or not through holes are formed in respective appropriate positions.\nThe inspection as to whether or not through holes are formed in respective appropriate positions has been normally carried out by a method in which a preform is observed from an end surface by use of an optical microscope.\nHowever, the conventional inspecting method by use of an optical microscope can inspect only forming positions of holes in the vicinity of an end surface of a preform. That is, the forming positions of the holes in the middle of end surfaces (at an arbitrary cross-section between both end surfaces) cannot be inspected. There is another method in which forming positions of through holes in the middle of a preform are estimated from (i) forming positions of holes in one end surface and (ii) forming positions of the holes in the other end surface. However, such estimation cannot accurately specify the forming positions of the holes in the middle of the end surfaces. In view of the circumstances, it is necessary to inspect which portion of a preform is suitable for manufacturing an optical fiber by cutting the preform in round slices in a case where forming positions of through holes in an end surface are largely shifted from respective appropriate positions. This inspection increases a manufacturing cost.\nOn the contrary, Patent Literature 5 discloses a method for inspecting whether or not through holes are formed in respective appropriate positions on the basis of an intensity distribution of forward scattered light which is generated by parallel light entering from a side surface of a preform."}
-{"text": "In general, military specifications define a fuze or fusing system as a physical system designed to sense a target or respond to one or more prescribed conditions such as elapsed time, pressure or command, and initiate a train of fire or detonation in munitions. Safety and arming are primary roles performed by a fuze to preclude ignition of the munitions before the desired position or time. For example, if a munitions is an array of energetic charges launched from a watercraft into shallow (or deep) water to clear mines and obstructions, the fuze should ideally determine that launch and deployment occurred as intended before initiating detonation. Several fuzes designed to just prevent in-air detonation are known in the prior art.\nU.S. Pat. Nos. 3,765,331 (the '331 patent) and 3,765,332 (the '332 patent) disclose water-armed air-safetied detonators in which a plurality of small explosives are aligned in a spaced-apart fashion in a fuze housing. The first of these small explosives is a delay charge which ignites when the ordnance is launched or released. The delay charge eventually burns and flashes down an adjacent housing bore to ignite a transfer charge. Detonation of the transfer charge releases energy that is either used to move a piston (the '331 patent) or is in the form of a shock wave (the '332 patent). This released energy is delivered to a chamber that is flooded with either air or water depending on the environment in which the fuze is immersed. Adjacent the flooded chamber is a firing pin/percussion primer (the '331 patent) or just a percussion primer (the '332 patent). If the flooded chamber is filled with air, the released energy in the form of a moving piston (the '331 patent) or shock wave (the '332 patent) will not transfer through to the next stage of the fuze. If, however, the fuze is submerged in water, the flooded chamber is filled with water and the released energy entering the flooded chamber is transferred therethrough to the next stage of the fuze.\nAlthough being air-safed, these devices still have several disadvantages. Use of stored energy (i.e., explosive material) for arming and firing is considered bad design practice because the energy is available at all times during storage and transportation, and may therefore be released due to unforeseen causes or situations. The use of explosives as part of the fuze train can be inherently problematic. These problems range from the safety concerns related to the construction and storage of such devices to the fact that these fuzes are not reusable.\nTo overcome the problems inherent with the use of explosives, a mechanical underwater firing mechanism is disclosed in U.S. Pat. No. 2,660,952. Briefly, a spring-loaded plunger is mounted in a housing. The head of the plunger is formed with a recess. Fitted in the housing coaxial with the plunger is a plug having a central bored portion in which a firing pin is temporarily positioned intermediately therein by a shear pin. As a result, small chambers are defined in the central bored portion on either side of the firing pin. The central bored portion of the plug opposes the plunger's recess and is sized at its exterior to fit within the recess. When the spring-loaded plunger is cocked, the head of the plunger is spaced apart from the central bored portion of the plug to define a chamber within the housing. An opening in the side of the housing at the chamber allows the environment surrounding the housing (e.g., air or water) to fill the chamber.\nWhen the spring-loaded plunger is released, the plunger recess envelops the central bored portion of the plug to compress any fluid trapped in the small chamber of the central bored portion between the firing pin and the head of the plunger. If the trapped fluid is air, the compression thereof will not develop forces sufficient to cause the shear pin to fail. However, if the trapped fluid is water, the compression forces imparted by the plunger will be sufficient to cause failure of the shear pin thereby allowing movement of the firing pin to impact a primer.\nWhile eliminating the use of explosives in the firing mechanism, this device has other disadvantages. For example, the requirement that a small chamber be defined in the central bored portion opposing the head of the plunger raises the possibility that an air bubble will form therein when the device is submerged in water. The presence of such an air bubble could prevent the mechanism from functioning underwater. At the same time, the requirement that the small chamber be present in the central bored portion could also bring about an unwanted firing. This could occur if the mechanism were not cocked and inadvertently dropped in water. Water could seep into the mechanism and fill the small chamber. Then, an in-air release of the (cocked) plunger could bring about movement of the firing pin just as if the mechanism were submerged in water. Another disadvantage brought about by the requirement of the small chamber arises in sub-freezing environments. Specifically, water in the small chamber could quickly freeze due to its small volume. If this occurs, the mechanism will not function."}
-{"text": "Zeolites are well known as typical examples of inorganic ion exchangers. Zeolites are capable of supporting various metals through cation exchange of Na+ ions or Ca.sup.2 + ions in the interior of their crystals for metal cations in an aqueous solution. Further zeolites have excellent characteristics with respect to specific surface area, heat resistance, water resistance, mechanical strength, etc. and are therefore widely used as gas adsorbing and separating agents, agents for treating waste water containing heavy metals, ion fixing agents and carriers for metal catalysts.\nZeolites are crystalline substances consisting primarily of SiO.sub.2 and Al.sub.2 O.sub.3, and the crystals thereof have three-dimensional structure and have regular minute pores. The composition is expressed generally by (M.sub.2, M')O. Al.sub.2 O.sub.3. mSiO.sub.2. nH.sub.2 O wherein M and M' are monovalent and bivalent metal ions, respectively, m is the coefficient of silica, and n is the coefficient of crystal water. Zeolites include the faujasite group (sodalite group), chabazite group and mordenite group which are different in the structure of crystals. The cations present in the voids or channels of zeolites can be exchanged for other metal ions, so that zeolites are utilized for softening hard water and separating off metal ions. When the alkali metal of zeolites is exchanged for bivalent or trivalent metal ions or hydrogen ions, the zeolite forms a strong solid acid, which is useful as an excellent catalyst for cracking petroleum and various carbonium ion reactions. Zeolites carrying silver, copper, zinc or the like supported thereon are kneaded with polymers to prepare antibacterial wrapping materials for preventing deterioration for use in the field of foods. Especially, the silver-incorporating zeolite is also excellent in ethylene-adsorbing ability and is therefore valuable for use.\nHowever, the zeolites incorporating silver, copper or zinc fail to fulfill the requirement of being inexpensive which is characteristic of inorganic materials because of material costs and complex manufacturing process, and are very expensive materials.\nAccordingly, the object to be achieved by the present invention is to develop a novel inorganic material which is usable in place of the silver-, copper- or zinc-incorporating zeolites and which has ethylene-adsorbing ability and antibacterial activity.\nThis object is fulfilled by exchanging Ca.sup.2 + ions, of these ions and monovalent cations, such as Na+ ions, which are present between the layers of a tobermorite for silver ions into the tobermorite although this has never been practiced in the past.\nA description will be given with reference to Na+ ion which is typical of monovalent cations.\nWe have directed attention to tobermorites which are fibrous layered compounds heretofore widely used as lightweight heat-insulating materials and noncombustible building materials because of their resistance to a high temperature of 650.degree. C. and excellent heat-insulating properties, and conceived the entirely novel idea that the Ca.sup.2 + ions, or these ions and Na+ ions which are present between the layers of the mineral can presumably be exchanged for silver ions by cation exchange. We have carried out intensive research based on this idea and consequently succeeded for the first time in the cation exchange of Ca.sup.2 + ions, or these ions and Na+ ions, which are present between the layers of a tobermorite, for silver ions. We have further found that this novel substance has ability to adsorb ethylene and antibacterial activity."}
-{"text": "\u201cData Warehouses\u201d are based on relational databases, and may include data management tools to extract data from various sources and manage that data. The data is transferred to the data warehouse in data packages that may include both data and metadata. Information related to the various data packages may be maintained in the data warehouse. All of the data and associated information is not necessarily held in the same database table. Rather, through the use of key fields, the data and associated information may be held in various tables that are linked together in an associated, relational database. The data warehouse operates in conjunction with, or \u201crides\u201d on, the database.\nProblems can arise when the links, that relate various data package information to each other, are broken. This can happen in a number of scenarios, some of which may not be avoidable, such as, for example, if a server shuts down at a particularly inappropriate moment. Often the reason for the link failure will not be known.\nA link failure does not necessarily cause an immediate problem with operation of the data warehouse. However, if an operation on, or with, a data package (\u201ca data request\u201d) is attempted after a link failure, the data request will likely fail. One such data request involves an attempt to delete a data package. Since the links are broken, the system, in some operations, does not \u201csee\u201d the data package. This will cause the system to overlook the data package and not delete it from, for example, an entity that holds the data for reporting (\u201ca data provider\u201d).\nProblems associated with such \u201cinvisible\u201d data packages, which can not be \u201cseen\u201d by the data warehouse system because of link failures, are call \u201czombie\u201d problems. One such zombie problem relates to requests involving such data. Such requests are called zombie requests. Zombie requests can cause major problems for the data provider because the request gets \u201clocked\u201d within the data warehouse or the data provider, is never fulfilled, and can not be deleted. However, most significantly, further data can not be loaded until the request is removed. Therefore this problem needs to be resolved with dispatch.\nOne way of addressing the problem of zombie requests is to have the database administrators remove all of the data and then reload the data. Another way of addressing the problem is to manually identify the data associated with the link failure. These techniques are labor intensive, time consuming, costly, and don't always succeed."}
-{"text": "The present invention pertains to a vehicle headrest, and particularly to one which is adjustable.\nThere exists a variety of proposed headrest constructions which allow a headrest to be adjusted to fit a particular driver or passenger. One adjustable headrest provides a pivoting adjustment as disclosed in U.S. Pat. No. 4,600,240 while a pneumatically operated bellows-type control is provided for another adjustable headrest as disclosed in U.S. Pat. No. 4,778,218. These headrests as well as others which include screw jack-type adjustments with scissor arms for supporting a movable plate with respect to a stationary plate typically are mounted to rods which extend into the seat back and provide vertical adjustment of the headrests.\nAlthough these systems provide the desired adjustment motion, some are either somewhat complicated to manufacture or require auxiliary equipment to operate and/or are relatively slow and cumbersome to use."}
-{"text": "In order to avoid heat losses through radiation, it is known to place a granular covering material on the top layer of a casting mold or ladle. The material used may be sawdust, glass, sodium silicate or Vermiculite. Recovered material such as aluminum oxide (Al2O3), calcium oxide (CaO) or aluminum drops may also be used.\nWhile the pouring of granular or powdered material onto the top layer of a casting ladle was initially carried out in a rather rough fashion by means of a hoisting mechanism and a chute, it was soon realized that it was necessary to find a more accurate solution. It is indeed desired to obtain accuracy in the metering of the desired quantity in order to obtain a predetermined layer thickness and also an even distribution of the bulk material over the whole surface to be covered.\nAn improved spreading device was suggested in EP 0 389 918, wherein the device comprises a fixed part from which a cover with scrapers is suspended, and a rotating part fitted below the scrapers. The rotating part comprises a series of slits that are uniformly distributed around a central shaft. The slits are radially widening in the direction away from the central shaft. Through the rotating movement of the rotating part and the levelling action of the fixed scrapers, the bulk material is uniformly distributed over the slits of the distribution bin. A rising and falling cover makes it possible to process varying volumes of bulk material.\nWhile this device achieves a homogeneous, protective and insulating covering layer on the top layer of a casting ladle containing molten steel it also has a number of disadvantages. Indeed, due to the central shaft, no material can be deposited directly to the center of the surface to be covered. This problem may however be partially solved by the use of deflectors for the bulk material.\nAnother disadvantage of the device is that after each operation, the distribution bin must again be filled with bulk material, which may cause some delay between two filling operations. Such a delay may e.g. cause unnecessary cooling of the molten steel or metal contained in the next ladle to be covered.\nFurthermore, it should also be noted that the construction of the device is rather cumbersome and necessitates a non-negligible amount of preventive and corrective maintenance work.\nAlthough, the device of EP 0 389 918 provides a homogeneous layer on the top layer of a casting ladle containing molten steel, there is still room for improvement."}
-{"text": "The invention is applied to a device for non-invasive measurement of heart rate information, in particular to a heart rate monitor used in connection with exercise and sports.\nThe measurement of heart beat frequency is an interesting field of application in connection with exercise. On the basis of the heart beat frequency, i.e. the heart rate, it is possible to obtain information e.g. on a persons stress level, recovery and development of physical condition, and consequently the proportion of training exercises and rest can be monitored and planned better.\nThe heart rate is measured on a person\"\"s skin on the basis of an electrocardiographic (ECG) signal generated by a heart beat. Additional information on ECG is available in the following publication by Guyton, Arthur, C., Human Physiology and Mechanisms of Disease, third edition, W. B. Saunders Company, 1982, ISBN 4-7557-0072-8, Chapter 13: The Electrocardiogram, which is incorporated herein as reference. The electrocardiographic signal is an electromagnetic signal originating from a heart beat, which is detected on the body of a person to be measured. The signal is measured by means of electrodes, which are in contact with the body at least at two points. By a polarization vector, the electrode that is located closest to the heart often acts in practice as the actual measuring electrode, while the other electrode serves as ground potential, to which the voltage measured by the measuring electrode is compared as a function of time.\nThe heart rate monitor electrodes to be placed on the chest are arranged in a known manner in a belt-like structure, i.e. a so-called electrode belt. The electrode belt is thus a ring-shaped attachment means that goes round the whole chest and can be tightened round the human chest. A structure of this kind is shown in FIG. 1. The electrode belts are known to have structures that comprise an electronic unit in the middle of the belt, with an electrode on both sides. The electrodes measure the electric pulse transmitted by the heart and transmit the measurement results to the electronic unit through an interface connecting the electrode and the electronic unit. The components included in the electrode belt, such as the electronic unit and the electrodes, are generally coated with plastic or rubber in order to protect the components against moisture, for instance. Depending on the structure of the electrode belt, the electronic unit often also comprises means for transmitting an electric pulse as an analog burst to a receiver and display unit worn on the wrist, for instance. Alternatively, the electrode belt itself may comprise the means for storing and displaying the electric pulses.\nIn general, the electrode belts have a structure in which the rubber or plastic support structure covering the components of the electrode belt is relatively rigid. These electrode belts are, in general, poorly suited for long-term, continuous use, and they chafe the skin easily. The belt-like structure of the electrode belt also limits its optimal positioning substantially at the heart with persons having large quantities of muscular or other tissue in the chest area. Also slim adults and children have troubles in wearing the rigid electrode belt, because it does not bend sufficiently to follow the contours of a human body with a narrow chest. In some prior art solutions, the problem is approached such that the plastic support structure between the electronic unit and the electrode is pleated, whereby the electrode belt bends immediately outside the electronic unit. However, this solution only reduces rigidity in bending the electrode belt, because the electrode belt is still ring-shaped and attachable round the chest.\nThe prior art solution has a serious drawback: it is difficult to fit the rigid electrode belt optimally round the chest to achieve the best measurement result, in particular in long-term, continuous use, when the electrode belt also chafes the skin easily.\nThe object of the invention is to provide an improved electrode structure and heart rate measuring arrangement for measuring an electrical heart beat signal on a human body such that the above problems can be solved. This is achieved with the following electrode structure for measuring an ECG signal on the chest of a person. The electrode structure comprises a band-like component having an inner surface to be placed against the skin of the person\"\"s chest and an outer surface opposite thereto, and which electrode structure comprises a first electrode at a first end and a second electrode at a second end of the electrode structure, and the inner surface of the electrode structure is an adhesive surface for attaching the electrode structure to the skin of the person\"\"s chest, and the electrode structure is arranged to measure a potential difference between the first and the second electrodes caused by the ECG signal.\nThe invention also relates to a heart rate measuring arrangement for measuring the ECG signal on the skin of a person\"\"s chest. The heart rate measuring arrangement comprises an electrode structure placed on a person\"\"s chest and a wrist-worn receiver unit, the electrode structure comprising a band-like component having an inner surface against the skin of the person\"\"s chest and an outer surface, opposite thereto, and which electrode structure comprises a first electrode at a first end and a second electrode at a second end of the electrode structure, the inner surface of the electrode structure being an adhesive surface for attaching the electrode structure to the skin of the person\"\"s chest, and the electrode structure being arranged to measure a potential difference between the first and the second electrodes caused by the ECG signal, the electrode structure further comprising ECG processing means communicating with the electrodes for measuring the potential difference between the first and the second electrodes caused by the ECG signal and for producing heart rate information on the basis of the measured potential difference, and the electrode structure further comprising a transmitter for transmitting the heart rate information to the wrist-worn receiver which comprises a receiver for receiving the heart rate information transmitted from the electrode structure, the wrist-worn receiver further comprising a display for presenting the heart rate information.\nThe preferred embodiments of the invention are disclosed in the dependent claims.\nIn the solution of the invention, it is intended that the electrode structure is placed on the skin of the user\"\"s chest. In one embodiment, the band-like component of the electrode structure is of flexible, soft material that fits the skin closely, and as a consequence it is comfortable and inconspicuous to wear and does not chafe the skin, and hence it is well suited for long-term use in ECG measuring. In one embodiment the band-like component of the electrode belt is continuous, in which both electrodes and their attachment means are integrated. The band-like component is disposable and economical to manufacture. In terms of design, it is a plaster-like sticker, for instance.\nThe electrode structure has a first electrode at a first end and a second electrode at a second end. The first and the second electrodes of the electrode structure are electrically separated from one another in order to enable the measurement of the potential difference between the electrodes. For optimal measurement of the heart rate signal, the first and the second electrodes should be located sufficiently far apart from one another so as to detect an electric ECG signal generated by a heart beat. The electrodes are thus advantageously placed at the ends of the electrode structure. Naturally, there may be more than said two electrodes.\nAccording to a preferred embodiment, the inner surface of the electrode structure is an adhesive surface for attaching the electrode structure on the skin of the person\"\"s chest. An advantageous manner to implement the adhesive surface and the electrodes is the embodiment, in which the electrodes located at the ends of the band-like component of the electrode structure are made of electrically conductive adhesive. Hence, the adhesive attaches the electrode structure on the skin of the person\"\"s chest. In a second embodiment the first electrode and the second electrode of the electrode structure consist of an electrically conductive membrane at both ends of the electrode structure, at the electrode. On the inner surface of the membrane, which is placed against the person\"\"s skin there is an electrically conductive adhesive. The adhesive is preferably an electrically conductive glue. Further, a third manner to implement the electrodes and the adhesive surface is an embodiment, in which the first electrode and the second electrode of the electrode structure consist of an electrically conductive membrane at both ends of the electrode structure, at the electrode, and the electrodes at both ends of the electrode structure are narrower than the band-like component. Around the electrodes, on the outer edges of the band, on the inner surface thereof, there is an adhesive, with which the electrode structure is attached to the person\"\"s skin. In this case, the adhesive need not be electrically conductive. Because the inner surface of the electrode structure is that portion of the electrode structure which is against the person\"\"s skin, the electrodes are preferably located on the inner surface of the band. The electrode structure can also be designed such that the electrodes are partly or completely located on both the inner surface and the outer surface of the band. Further, the electrode structure can be designed such that the electrodes are located on the inner surface of the band, but they have interfaces also on the outer surface of the band.\nThe electrode structure also comprises an electronic unit communicating with the electrodes. The electronic unit is an electronic component attached to the band-like part of the electrode structure with one or more gripping means. The electrodes of the electrode structure communicate with the electronic unit. The electronic unit comprises ECG processing means, by which the potential difference caused by the ECG signal between the first and the second electrodes is measured, and an estimate for detected heart beat time instants is formed from the heart rate signals measured by the electrodes, and further the heart beat rate is calculated on the basis of the detected heart beat time instants. The electronic unit also comprises a transmitter for transmitting heart rate information to a wrist-worn receiver, which comprises a receiver for receiving the heart rate information transmitted from the electrode structure, and a display for presenting the heart rate information.\nThe wrist-worn receiver is located in a watch-like device that the user wears on his wrist, such as a heart rate monitor or a wrist computer. Transmission of information between the electrode structure and the heart rate monitor is thus carried out in known manners, for example through a connecting line, optically or electromagnetically. In the embodiment in question, the display for presenting the heart rate information is also preferably located in the wrist-worn receiver.\nThe electronic unit is preferably arranged in a casing which comprises one or more gripping means for attaching the electronic unit to the stap-like component of the electrode structure. The gripping means are most preferably located on that surface of the electronic unit casing which is against the person\"\"s skin. The preferable structures of the gripping means include attachment slots, a pivoted, clamping clip, or the like, that are in the central unit casing. The gripping means or the central unit casing comprises conductive connecting means, through which the ECG signal measured with the electrodes is applied from the electrodes to the electronic unit.\nThe invention also has an advantage that the electrode structure is inconspicuous, comfortable and well suited for long-term use as compared with the known solutions."}
-{"text": "Depression has been a major social problem. A survey in the United States of America estimates that the lifetime prevalence rate of depressive disorders is 26% for women and 12% for men (Non-patent Document 1). An extensive survey was also carried out in Japan, and it was reported that the lifetime prevalence rate was 14%. It was also reported that 84% of the patients with depression have complication, 61% of the patients have other mental diseases, 30% of the patients have personality disorders, and 58% of the patients have complication of physical diseases (Non-patent Document 2).\nCommon symptoms of depression include a wide variety of symptoms, such as depressive mood, lethargy, depressive attitudes, suicidal ideation, impatience, insomnia, anorexia, decreased sexual impulse, and physical ailments. Pathologically-severe depression is called major depression. With regard to major depression, there is reported experimental data showing that in many cases, a rapid recovery is observed during the first 40 weeks, but the recovery plateaus thereafter. It is reported that results of a 20-year follow-up survey on patients having major depression show that 15 to 19% of the patients had chronic residual symptoms such as working difficulty (Non-patent Document 3).\nWith regard to causes of depression, various theories are suggested. A hypothesis that was suggested in the early days and is still popular is a theory stating that a decrease in the function of monoamines in the brain is responsible for depression (this hypothesis is called \u201cintracerebral monoamine hypothesis\u201d). This theory is suggested on the basis of the fact that a tricyclic antidepressive agent having an action to enhance the function of monoamines in the brain by blocking the reuptake of monoamines in cells is effective for remission of depression.\nCurrently, especially the function of serotonin, which is an intracerebral monoamine, attracts great attention, and there is a popular hypothesis suggesting that the antidepressive action can be produced by sensitizing serotonin-1 receptors present in the serotonin nerve cell body to increase the amount of free serotonin from presynaptic cells (Non-patent Document 4). This idea is basically the same as that of the classical intracerebral monoamine hypothesis. In any cases, it is considered that depression is caused by a decrease in the intracerebral serotonin function.\nWith regard to medications for depression, first, a tricyclic antidepressive agent, imipramine, and monoamine oxidase inhibitors were developed about 40 years ago. On the basis of the idea that the action to increase serotonin and noradrenaline in the synaptic cleft, which is one of the features the above drugs, is the core of the antidepressant effect, many antidepressive agents have been developed following imipramine and monoamine oxidase inhibitors. In recent years, use of selective serotonin reuptake inhibitors (SSRI) and serotonin/norepinephrine reuptake inhibitors (SNRI), which more selectively inhibit the reuptake of the monoamines by serotonin transporters, have been used in clinical practice. There are as many as 17 antidepressive agents that are approved in Japan at this time, including tricyclic/tetracyclic antidepressive agents, SSRI, and SNRI (Non-patent Document 5).\nThese antidepressive agents are effective against various diseases besides depression. According to package inserts of the pharmaceutical agents presented by the Pharmaceuticals and Medical Devices Agency, for example, fluvoxamine maleate, which is an SSRI, is effective against obsessive-compulsive and social anxiety disorders; paroxetine hydrochloride, which is an SSRI, is effective against panic and obsessive-compulsive disorders; and sertraline hydrochloride, which is an SSRI, is effective against a panic disorder. Further, amitriptyline hydrochloride, which is a tricyclic SNRI, is effective against nocturnal enuresis.\nHowever, it is reported that these antidepressive agents have side effects. For example, nausea, headache, hypersensitivity and the like are reported as side effects of SSRI, and tremor, tachycardia, erection/ejaculation disorders and the like are reported as side effects of SNRI. There are several reports that side effects are increased in the case of administration of multiple antidepressive agents, compared with administration of a single antidepressive agent (Non-patent Document 6).\nIt is also reported that when administration of an antidepressive agent is discontinued, withdrawal symptoms may occur. Withdrawal symptoms of tricyclic/tetracyclic antidepressive agents that are reported include: digestive symptoms and physical discomfort accompanied by anxiety and irritability, such as lethargy, emesis, and headache; sleep disorders such as insomnia and multi-dream; motility disorders such as akathisia and Parkinsonian symptoms; activation of behavior that leads to a transition to a manic state; and arrhythmia. These symptoms, except for arrhythmia, are also observed in withdrawal symptoms of SSRI (Non-patent Document 7). Disturbance to the sense of equilibrium, paresthesia, aggressive/impulsive behavior and the like are reported as withdrawal symptoms characteristic of SSRI (Non-patent Document 8).\nAs described above, it is reported that the antidepressive agents that are widely used at this time have side effects, and the development of safe food ingredients and components, which have ever been eaten and are applicable in place of the above antidepressive agents, are expected. Royal jelly-derived compositions (Patent Document 1) and compositions comprising a hop extract as a main component (Patent Document 2) are reported as food-derived compositions having antidepressive action.\n1-Methyl-1,2,3,4-tetrahydro-\u03b2-carboline-3-carboxylic acid (MTCA) is a compound contained in fruits, such as lemon, grapefruit, orange, and mandarin (Non-patent Document 9), beer, wine (Non-patent Document 10), soy sauce (Non-patent Document 11) and the like, and two types of diastereomers, (1S,3S) and (1R,3S), are known. It is known that MTCA is produced nonfermentatively in the presence of tryptophan and acetaldehyde; especially, it is known that (1S,3S)-MTCA (also called (1S,3S)-2,3,4,9-tetrahydro-1-methyl-1H-pyrido[3,4-b]-indole-3-carboxylic acid) is produced through fermentation by S. cerevisiae (Non-patent Document 12). Further, it is known that (1S,3S)-MTCA has antioxidation action, prevents self-polymerization reaction of \u03b3-crystallin by the Fenton reaction using FeCl3, and prevents photopolymerization reaction of \u03b3-crystallin (Non-patent Document 13). Further, there is a description that pinoline (6-methoxy-1,2,3,4-tetrahydro-\u03b2-carboline), which is a carboline compound, inhibits the activity of monoamine oxidase A, and also inhibits the uptake of serotonin in the brain (Non-patent Document 14). Further, it is reported that pinoline reduces immobility time in forced-swimming tests using rats, and that this is the same action as that of a tricyclic antidepressive agent (Non-patent Document 15).\nHowever, inhibitory action of serotonin uptake, antidepressive effect, and improvement of learning motivation by MTCA have completely been unknown."}
-{"text": "1. Field of the Invention\nThe present invention relates to a method and apparatus for processing a linework image, and more particularly to a method and apparatus for finding and reconnecting a break in a linework image.\n2. Description of the Related Art\nA block copy is prepared as an original of characters and lineworks in a prepress process for a color printing plate. The block copy is formed by arranging phototype-setting characters and drafted keylines on a layout sheet in the same dimensions and quality as the final products. Instructions for the subsequent processes are also given on the layout sheet; that is, the layout sheet acts as an instruction paper for the prepress process.\nTint laying is generally performed in the prepress process. Tint laying is processing in which a specified region within an image is uniformly filled with a desired color. Automatic tint laying is employed in some of the modern image processing systems. In such systems, tint laying is completed, for example, by obtaining binary image data of a block copy image with an image scanner and filling a certain closed region in the block copy image with a desired color with an image processor.\nA linework, which is supposed to define a closed region in a block copy image, is sometimes found incomplete or broken when filling of the closed region is executed. The incomplete or broken linework is generally attributable to an incomplete draft or layout sheet, or unsuccessful image reading with an image scanner.\nLine breaks are generally reconnected manually; an operator directly detects line breaks in a block copy image displayed on a CRT and fills them to correct the block copy image.\nA line break has generally a width of one to several pixels. Accordingly, only experienced operators can efficiently detect breaks on the CRT. When there are several breaks in a linework, it takes rather a long time to detect and reconnect all of them.\nJAPANESE PATENT LAYING-OPEN GAZETTE No. Sho-61-139892 proposes a method of reconnecting line breaks which utilizes a technique of changing binary image data into vector data. Vector data processors for executing this method are rather complicated and expensive.\nOn the other hand, those image processors for processing lineworks (hereinafter referred to as linework processors), which change the level of binary image data for each pixel, are relatively simple. Accordingly, addition of a specific function of vector data processing to linework processors drastically increases the price of the processors.\nFurthermore, the method including vector data processing does not perfectly detect line breaks and requires operators confirmation for the existence thereof. It means that the total time for processing is not shortened greatly."}
-{"text": "The present invention relates to illumination devices generally, and more particularly to portable illumination devices for removable attachment to various objects. This invention also relates to methods for removably attaching portable illumination devices to various objects.\nA variety of illumination devices have been implemented to illuminate areas which receive no or insufficient light for human utilization. Most commonly, artificial means have been used to illuminate such areas. Some of the devices that have been developed to provide such portable illumination include electric or oil lanterns, self-supporting lights, and individual lighting units powered by batteries or electrical cords.\nA particular application for such illumination means is in illuminating desired areas for inspection and/or manipulation, where traditional lighting systems are not available. Relatively compact illumination devices, therefore, have become desirable for such purposes.\nFlashlights, in particular, have become very popular in such situations for their relative portability and low cost. One of the most common applications for such flashlights is in typical household maintenance and construction, though a wide variety of other flashlight applications are used on a daily basis. Often times, the flashlight operator utilizes the flashlight to illuminate an area in which tools or other implements are being utilized. To effectively manipulate such tools, the operator has traditionally needed to position the flashlight such that the flashlight illuminates the work area without requiring the user to continuously hold the flashlight. However, most portable flashlights project a relatively narrow illumination beam, thereby forcing the user to repeatedly re-direct and re-position the flashlight to illuminate desired areas. Thus, a need arose for a means of directing an illumination beam from a flashlight without requiring continuous manipulation from the user\"\"s hands.\nA number of mechanisms have been introduced which are designed to grasp conventional flashlights, and to be placed in various positions to more effectively illuminate desired areas. Examples of such devices are described in U.S. Pat. Nos. 4,214,688, 2,354,853, and 1,540,372. Such devices provide means for gripping conventional flashlights such that user manipulation is not required, thereby freeing the users hands to perform other tasks. However, such devices typically have limited uses, and may not be universally adapted to function in various applications.\nPortable illumination mechanisms proposed to date have not addressed the issue of having a lightweight device that may be removably attached to various objects such that the object may still be used for their intended purpose while having such an illumination device attached thereto. Furthermore, no such device has been proposed which incorporates an attachment means in the body of the portable illumination device itself, so that a single illumination device may be easily transferred to and from several receptacle objects.\nAccordingly, it is a principle object of the present invention to provide a portable illumination device having attachment means integrated therein for removable attachment of the illumination device to various objects.\nA further object of the present invention is to provide a portable illumination device that may be easily secured to and removed from small objects without interfering with the intended functionality of such objects.\nA still further object of the present invention is to provide a compact, portable illumination device that is removably attachable to various hand tools, including tape measures.\nA yet further object of the present invention is to provide a portable illumination device having universal attachment means integrated therein.\nAnother object of the present invention is to provide a portable illumination device which incorporates a plurality of distinct illumination means.\nA still further object of the present invention is to provide a portable illumination device which incorporates multi-directional illumination means, each of which illumination means has a distinct illuminatory purpose.\nA yet further object of the present invention is to provide a portable illumination device which incorporates a plurality of illumination means each emitting designated illumination wavelengths.\nA still further object of the present invention is to provide a portable illumination device incorporating one or more light emitting diodes (LED).\nBy means of the present invention, portable illumination is enhanced by providing an illumination device having attachment means integrated therein for removable attachment of the illumination device to various objects. Furthermore, such a portable illumination device may be utilized in a manner so as to avoid interfering with intended functionality of respective objects that the device is removably attached to. Such objects include everyday tools such as tape measures, saws, drills, nail guns, screw guns, as well as other common objects such as toolboxes, clothing belts, walls, and building studs. The portable illumination device of the present invention preferably utilizes one or more types of illumination elements, including light emitting diodes, incandescent lights, fluorescent lights, or other gas-charged lighting elements. Such illumination elements are preferably positioned in the illumination device and directed to illuminate specific target areas.\nOne embodiment of the portable illumination device of the present invention includes an outer frame having illumination means disposed at least partially therein, and attachment means for removably attaching the illumination device to desired objects, with the attachment means being formed integral with the outer frame. The device is preferably sized and configured to receive specific securing means which are connected to the desired objects. The attachment means may include one or more slots which are sized and configured to receive securing means being made up of protrusions sized and configured accordingly. The illumination means of the device may include one or more individual light emitting diodes (LED), wherein the light emitting diodes may illuminate in one or more distinct colors. The device may be particularly sized and configured to be removably attachable to a side surface of a tape measure such that the illumination device does not interfere with the functionality of the tape measure.\nThe present invention also contemplates a method of illuminating a target zone. The method includes providing a portable illumination device having attachment means formed integral therewith, and a securing means that is sized and configured to receive the attachment means. The method further includes connecting the securing means to a desired object, removably attaching the illumination device to the securing means via the attachment means, and actuating the illumination means. The illumination device of the present invention may also be utilized independent of the securing means, such that the illumination device is actuatable and directable in a variety of hand-held applications."}
-{"text": "As more women become aware that breastfeeding is the best source of nutrition for a baby, and also offers health benefits to the nursing mother, the need is increasing for breast pump solutions that are user-friendly, quiet, discrete and versatile for use by a nursing mother in various situations. This is particularly true for the working mother, who is away from the home for eight to ten hours or more and needs to pump breast milk in order to have it available for her baby, but it is also a requirement for many other situations where the mother is away from the privacy of the home for an extended period, such as during shopping, going out to dinner or other activities.\nAlthough a variety of breast pumps are available, most are awkward and cumbersome, requiring many parts and assemblies and being difficult to transport. Hand pump varieties that are manually driven are onerous to use and can be painful to use. Some powered breast pumps require an AC power source to plug into during use. Some systems are battery driven, but draw down the battery power fairly rapidly as the motorized pump continuously operates to maintain suction during the milk extraction process. Many of the breast pumps available are clearly visible to an observer when the mother is using it, and many also expose the breast of the mother during use.\nThere is a continuing need for a small, portable, self-powered, energy efficient, wearable breast pump system that is easy to use and is discrete by not exposing the breast of the user and being invisible or nearly unnoticeable when worn."}
-{"text": "The present invention relates to a plug part for establishing an optical and/or electrical plug-in connection with a mating plug part by plugging together the plug part with the mating plug part in at least one plugging-together direction, the plug part having a plug part housing and at least one optical and/or electrical contact and an elastic opening prestressing element, acting in a direction parallel to the plugging-together direction, and an opening element for opening a flap on the mating plug part, the opening element being mounted displaceably in relation to the plug part housing in directions parallel to the plugging-together direction in or on the plug part housing, and being prestressed in one of these directions parallel to the plugging-together direction by the opening prestressing element.\nMany configurations of plug parts for establishing optical and/or electrical plug-in connections are known. Not only optical and/or electrical signals but also power can be transmitted by way of the optical and/or electrical plug-in connections. In the case of more sensitive contacts, it is already known in the prior art to protect them from external influences in the unplugged state of the plug part and the mating plug part by a flap. As an example, reference is made here to EP 2 146 232 A1. In this document, both the plug part and the mating plug part each have a flap protecting the contacts, which is only opened during the plugging together of the plug part and the mating plug part."}
-{"text": "Field\nEmbodiments relate generally to sodium free glasses and more particularly to fusion formable sodium free glasses which may be useful in photochromic, electrochromic, Organic Light Emitting Diode (OLED) lighting, or photovoltaic applications, for example, thin film photovoltaics.\nTechnical Background\nThe fusion forming process typically produces flat glass with optimal surface and geometric characteristics useful for many electronics applications, for instance, substrates used in electronics applications, for example, display glass for LCD televisions.\nOver the last 10 years, Corning fusion glass products include 1737F\u2122, 1737G\u2122, Eagle2000F\u2122, EagleXG\u2122, Jade\u2122, and Codes 1317 and 2317 (Gorilla Glass\u2122). Efficient melting is generally believed to occur at a temperature corresponding to a melt viscosity of about 200 poise (p). These glasses share in common 200p temperatures in excess of 1600\u00b0 C., which can translate to accelerated tank and electrode corrosion, greater challenges for fining due to still more elevated finer temperatures, and/or reduced platinum system life time, particularly around the finer. Many have temperatures at 3000 poise in excess of about 1300\u00b0 C., and since this is a typical viscosity for an optical stirrer, the high temperatures at this viscosity can translate to excessive stirrer wear and elevated levels of platinum defects in the body of the glass.\nMany of the above described glasses have delivery temperatures in excess of 1200\u00b0 C., and this can contribute to creep of isopipe refractory materials, particularly for large sheet sizes.\nThese attributes combine so as to limit flow (because of slow melt rates), to accelerate asset deterioration, to force rebuilds on timescales much shorter than product lifetimes, to force unacceptable (arsenic), expensive (capsule) or unwieldy (vacuum fining) solutions to defect elimination, and thus contribute in significant ways to the cost of manufacturing glass.\nIn applications in which rather thick, comparatively low-cost glass with less extreme properties is required, these glasses are not only overkill, but prohibitively expensive to manufacture. This is particularly true when the competitive materials are made by the float process, a very good process for producing low cost glass with rather conventional properties. In applications that are cost sensitive, such as large-area photovoltaic panels and OLED lighting, this cost differential is so large as to make the price point of LCD-type glasses unacceptable.\nTo reduce such costs, it is advantageous to drive down the largest overall contributors (outside of finishing), and many of these track directly with the temperatures used in the melting and forming process. Therefore, there is a need for a glass that melts at a lower temperature than those aforementioned glasses.\nFurther, it would be advantageous to have a glass useful for low temperature applications, for instance, photovoltaic and OLED light applications. Further, it would be advantageous to have a glass whose processing temperatures were low enough that the manufacturing of the glass would not excessively consume the energy that these applications are aiming to save."}
-{"text": "This invention relates to treatment of catalyst supports.\nSupported chromium oxide catalysts have become an important vehicle for polymerization of 1-olefins such as ethylene and predominantly ethylene copolymers. As originally commercialized, these polymerizations were carried out under solution conditions. However, it was early recognized that a more economical route to such polymers was to utilize a slurry system sometimes called a particle form system, wherein the polymerization is carried out at a low enough temperature that the resulting polymer is insoluble in the diluent. A great deal of technology has been developed on means to adapt catalysts to compensate for the inherent difficulty in slurry polymerization systems of achieving high enough flow in the resulting polymer. However, another polymer property which is important in many applications is the molecular weight distribution which has an effect on the response to shear. There are a substantial number of applications which require a higher shear response polymer than is normally produced with chromium catalysts in a slurry polymerization system."}
-{"text": "The present invention relates to an optical scanner, and particularly to the optical scanner by which the light is 2-dimensionally scanned.\nPresently, the optical scanner by which the ray such as the laser light is deflected and scanned, is used for an optical device such as a bar-code reader, laser printer, display. Conventionally, as the optical scanner such as this, there is a polygonal mirror in which a polygonal column mirror is rotated by a motor, and a galvano-mirror in which a plane mirror is rotation-oscillated by an electromagnetic actuator. However, in a mechanical structure which is driven by the motor or electromagnetic actuator, because the shape of the structural parts is large or cost is high, there is a problem that the downsizing of the device for which the optical scanner is used, is hindered, or results in high cost. Further, when the ray is 2-dimensionally scanned, generally, a device in which the polygonal mirror and the galvano-mirror are combined is used, however, in order to conduct the accurate 2-dimensional scanning, it is necessary that the respective scanning directions are accurately positioned so that the scanning directions by respective mirrors are orthogonal with each other, and there is also a problem that the optical adjustment is very complicated.\nTherefore, in order to realize the downsizing, low cost, and the increase of the productivity, by using a micro-machining technology by which silicon or glass for which the semiconductor manufacturing technology is applied, is micro-machined, various micro optical scanners in which the structural parts such as the mirror or an elastic beam are integrally formed on the semiconductor substrate are successively developed.\nFor example, two frames of an inner frame for supporting the mirror through the first double holding beam, and an outer frame for supporting the inner frame through the second double holding beam are provided, the first and the second double holding beams are orthogonal with each other, and the bending-ness of the second double holding beam and the resonance frequency of the oscillation of the first double holding beam, are made close. Then, the magnetic distortion film is formed on the second double holding beam, and the engineering of the optical scanner (refer to Japanese Patent Unexamined Publication: Tokkai No. 2003-207737) by which, when the distortion and the bending oscillation are generated, the mirror is rotated by using the first and second double holding beams as 2 axes, and 2-dimensionally optical scanning is conducted, further, in almost the same structure] as the optical scanner disclosed in Patent Document 1, 2 frames and 2 double holding beams are provided, and by a plurality of piezoelectric elements arranged in the inner frame and the outer frame, when the rotation torque is respectively independently, acted on the inner frame and the mirror, the engineering of the optical scanner (refer to Tokkai-No. 2005-148459) which rotate the mirror when the first and the second double holding beams are made 2 axes, is disclosed.\nHowever, the optical scanner disclosed in Patent document 1, the distortion and bending oscillation are simultaneously conducted on the second double holding beam by the actuator.\nThat is, because distortion and bending can not be independently conducted, it is considered to be difficult that the horizontal and vertical deflection are independently controlled. Further, because the resonance frequencies of the distortion oscillation of the first double holding beam and the bending oscillation of the second double holding beam are set so that they are almost same value, when the horizontal and the vertical deflection frequencies are largely different, that is, when the resonance frequencies of the distortion and the bending oscillation are largely different, it is considered to be difficult that the first double holding beam and the second double holding beam are resonated at respectively necessary deflection frequency. Further, the optical scanner disclosed in Patent Document 2, because, as for the horizontal deflection and the vertical deflection, a plurality of piezoelectric element are respectively independently provided, there is a possibility that it results in the complexity of the device, and high cost."}
-{"text": "Field of the Invention\nThe invention relates to a method and system for producing screen or raster data for imaging units of a printing machine.\nPrinting machines with imaging units have become known heretofore. Imaging units of this type produce a printing image on a printing plate of the printing machine by scanning with a laser beam, the position and size of the screen or raster dots of the printing image being controlled by screen or raster data, with which the laser beam is modulated in the course of a scanning movement over the surface of the printing plate.\nWith the aid of imaging units of this type, as compared with the conventional film exposure technique, a plurality of processing stages on the path from an original image to the finished print are dispensed with, so that alternating print jobs can be processed speedily and cost-effectively. Nevertheless, even when such imaging units are used, the processing process from the original image to the finished print breaks down into a multiplicity of processing stages which have to be performed successively and which, to some extent, are time-consuming. In a conventional screen production system, as illustrated in FIG. 1 of the drawings herein, raw image data initially passes through a screen processor 1 (screen image processor RIP), which converts this data into a plurality of partial images, respectively, one partial image for each color to be printed, while simultaneously performing all of the necessary calibration operations. These partial images, respectively, are transferred by an administration unit 2 into a buffer storage unit 3. After the screen processor 1 has completely processed the original image, and all the necessary partial images are present in the buffer storage unit 3, they can be transferred therefrom to the imaging units 4 of a diagrammatically illustrated printing machine 5. Buffering the screen data in the unit 3 is necessary, because the imaging units 4, respectively, must be supplied with data without interruption at a data rate which is predefined or prescribed by the type of construction thereof, and a non-buffered imaging process will fail if the screen processor 1 is not capable of supplying the data at the required rate.\nConsequently, the imaging units 4 remain inactive while the screen processor 1 is operating, and the reverse. Because the screen data for each screen dot to be produced by the imaging units must specify the size of the respective dot, they include a significantly greater quantity of data than the raw data originally input into the screen processor 1, which, for example, can be a file in Postscript or PDF format. Storing the screen data for the situation wherein they are used again for a subsequent print job is therefore rather complicated. In addition, if the screen processor 1 performs a high-quality calibration, the dot increases or growths of the individual imaging units of the various stages of the printing machine are taken into account in the specification of the screen dots. These dot increases or growths can vary from one imaging unit 4 to another, within certain limits, because of the scatter in the optical properties of the laser and the scanning system, which are used. Consequently, the screen processor 1 produces the screen data for a given printing ink, respectively, taking specifically into account the properties of the imaging unit 4 which sets the image on the printing plate for the relevant color. If a print job is to be repeated based upon stored screen data, it is consequently necessary for the distribution of the printing inks to the various stages of the printing machine 5, at the time the print job is repeated, to be the same as that for which the screen data were originally calculated. If, in the interim, jobs with a different sequence of colors have been processed, this would require lengthy washing and re-inking of the printing machine.\nFor a repetition of the print job on a different machine, renewed screening calibrated to the imaging units of the machine is always necessary, even if the machines are of identical construction and are equivalent in all the other relevant parameters, such as the printing material, the color used, the plate material, and so forth.\nIt is accordingly an object of the invention to provide a method and system for producing screen data for imaging units of a printing machine, which permit further acceleration of operating steps from the original raw image to the print.\nWith the foregoing and other objects in view, there is provided, in accordance with one aspect of the invention, a method of producing screen data for imaging units of a printing machine, which comprises the steps of breaking down raw image data into a plurality of partial images, respectively, corresponding to one printing ink; outputting the partial images to a plurality of screen processors, respectively, corresponding to the number of colors to the printed; and processing the partial images simultaneously for forming screen data by the screen processors for output to a respective one of the imaging units.\nIn accordance with another mode, the method of the invention includes performing by the screen processors, respectively, a calibration of the partial images relative to dot growth of the imaging unit with which it is associated.\nIn accordance with a further mode, the method of the invention includes breaking down the raw image data in a preprocessing unit that is separate from the screen processors.\nIn accordance with an added mode, the method of the invention includes performing a calibration by the preprocessing unit, relative to properties of at least one of the printing material, the printing ink, the blanket cylinder, and the plate material of the printing machine.\nIn accordance with an additional mode, the method of the invention includes, by the preprocessing unit, outputting partial image data to the screen processors before the step of breaking down the raw image data has been finished.\nIn accordance with yet another mode, the method of the invention includes buffering the screen data before outputting the screen data to the respective imaging units.\nIn accordance with yet a further mode of the method invention, the step of breaking down raw image data includes a trapping procedure.\nIn accordance with a concomitant aspect of the invention, there is provided a screen production system for a printing machine having imaging units, comprising a preprocessing unit for receiving raw image data and for breaking down the raw image data into a plurality of partial images, respectively, corresponding to one printing ink, and a plurality of screen processors, respectively, for screening the partial images, respectively, and for outputting to the imaging units, respectively, partial-image screen data obtained from the respective partial images.\nThus, instead of determining, in order, the sizes of the screen dots of the colors to be printed for each point or dot in an image to be printed here, initially, a separation into various colors is performed, and the partial image data, respectively, corresponding to one color, is processed by a screen processor which is specifically provided for the purpose and which is therefore capable of providing the screen data of the associated color significantly more rapidly than heretofore.\nThe calibration of the partial images relative to the dot growth of the imaging unit, the calibration being necessary for a high-quality print, is preferably performed by the screen processor which is associated with the imaging unit. Those partial images not yet screened, which form data files of a relatively small size, can comfortably be stored intermediately for a repetition of the print job. When the print job is to be repeated, the intermediately stored partial images are screened anew. If, in the interim, the distribution of the colors to the stages of the printing machine has changed, or another printing machine is to be used, the renewed screening can be performed by a screen processor that differs from the first one and can be calibrated to the dot growth of a different printing stage.\nThe step of breaking down the raw image data into partial images is preferably performed in a preprocessing unit that is separate from the screen processors. This preprocessing unit can perform a calibration relative to all those printing parameters which are standard for the various stages of the printing machine, thus, for example, the printing material used, the printing ink, the blanket cylinder or the plate material of the printing machine.\nA particular advantage in using the preprocessing unit is the possibility of intermeshing in time the actions of breaking down and screening, by the preprocessing unit outputting partial image data to the screen processors before the step of breaking down the raw image data has been finished completely.\nIn order to be able to ensure a uniform data stream at the rate required by the imaging unit, the screen data are preferably buffered before being output to the imaging unit. However, such a buffer can be considerably smaller than that needed to store an entire partial image. If the computing power of the screen processors is sufficiently high so that they can ensure a stream of screen data at the rate required by the imaging unit, the buffer can even be omitted completely.\nThe step of breaking down the raw image data preferably also includes a trapping procedure.\nThe object of the invention is also achieved by a screen production system for a printing machine which includes a preprocessing unit for receiving raw image data and for breaking down the raw image data into a plurality of partial images, each of which corresponds to one printing ink, and a plurality of screen processors for screening one of the partial images, respectively, and for outputting the thus obtained partial-image screen data to an imaging unit.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as a method and embodied as a system for producing screen data, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawings, wherein:"}
-{"text": "The present invention relates to surgical instruments and, in various arrangements, to surgical stapling and cutting instruments and staple cartridges for use therewith that are designed to staple and cut tissue.\nCorresponding reference characters indicate corresponding parts throughout the several views. The exemplifications set out herein illustrate various embodiments of the invention, in one form, and such exemplifications are not to be construed as limiting the scope of the invention in any manner."}
-{"text": "Although the present invention can be applied to a wide range of plastics material parts for decorative and trim purposes, for example in the field of automotive engineering, but also in other technical fields, for example when equipping any vehicles, aircraft or buildings, and also in domestic appliances or entertainment devices, the invention and the problem on which it is based are described in detail below with reference to a trim component for a motor vehicle, but without restricting the invention to that effect.\nIn the context of motor vehicles, in particular cars, vendors and users now require not only efficient drive technology, low fuel consumption and the like, but also a high-grade and aesthetic appearance. Decorative or trim components for the interior of motor vehicles are currently often produced from plastics material, at least in part. This allows the cost-effective production of complicated geometries with a simultaneous reduction in weight.\nIn this context, it can happen that a trim component of this type is of a high-gloss design for aesthetic reasons and is to be embellished in particular, predefined regions with a decoration, for example a chrome-look decoration.\nFor this purpose, the trim component can be implemented in two or more separate pieces, one of the pieces being provided with a high-gloss surface and a second of the pieces being provided with a chrome finish or another decoration. For this purpose, the second piece may for example be chromium-plated. However, a two-piece or multi-piece design is associated with high complexity and high costs, especially since the pieces must be produced individually and then assembled.\nThis is a situation which is worth improving."}
-{"text": "Display screens of electronic devices such as a mobile phone, personal data assistants (PDAs), personal computers, game consoles and so on, are developing towards big sizes. So, the electronic device generally includes a user interface with a big size."}
-{"text": "Today's hospitality systems usually support a subscription of video or audio information like movies, music etc. which are displayed by an in-room screen or played by in-room loudspeakers.\nA disadvantage of current prior art hospitality systems is the limited way to deal with preferences of individual guests regarding the lighting of the guest's room, temperature etc.\nYet another disadvantage of known hospitality systems is that an identical selection of information and entertainment contents is offered to all guests. Every guest is encountered by a bulk of selectable entertainment media and other contents which will limit the usability of selecting media in the guest's sphere of interest.\nA further disadvantage of known hospitality systems is the limited manner to report alarm situation to guests and to effectively guide personnel and guests in the case of an evacuation."}
-{"text": "1. Field of the Invention\nThis invention relates to an improved method for preparation of certain zeolites having a high silica/alumina mole ratio, i.e. greater than about 10. The improvement resides in adding a source of aluminum ions to a silica-rich amorphous reaction medium at a slow and controlled rate whereby the concentration of aluminum ions in the amorphous phase of the reaction mixture is maintained at steady state during crystallization.\nMore particularly, this invention relates to the above improved method of preparation of certain zeolites whereby crystallization time is substantially reduced from that required when conventional prior art methods of preparation are utilized and the resulting zeolite exhibits somewhat improved steam stability, i.e. stability toward steam deactivation and dealuminization.\nEven more particularly, this invention relates to organic compound conversion, such as, for example, catalytic conversion of oxygenates, such as methanol, and syn-gas conversion where water is always present during reaction, with the improved zeolite product of the present improved method as a catalyst.\n2. Discussion of Prior Art\nZeolitic materials, both natural and synthetic, have been demonstrated in the past to have catalytic properties for various types of hydrocarbon conversions. Certain zeolitic materials are ordered, porous crystalline aluminosilicates having a definite crystalline structure within which there are a large number of smaller cavities which may be interconnected by a number of still smaller channels. Since the dimensions of these pores are such as to accept for adsorption molecules of certain dimensions while rejecting those of larger dimensions, these materials have come to be known as \"molecular sieves\" and are utilized in a variety of ways to take advantage of these properties.\nSuch molecular sieves, both natural and synthetic, include a wide variety of positive ion-containing crystalline aluminosilicates. These aluminosilicates can be described as a rigid three-dimensional framework of SiO.sub.4 and AlO.sub.4 in which the tetrahedra are cross-linked by the sharing of oxygen atoms whereby the ratio of the total aluminum and silicon atoms to oxygen is 1:2. The electrovalence of the tetrahedra containing aluminum is balanced by the inclusion in the crystal of a cation, for example, an alkali metal or an alkaline earth metal cation. This can be expressed wherein the ratio of aluminum to the number of various cations, such as (Ca/2), (Sr/2), Na, K or Li is equal to unity. One type of cation may be exchanged either entirely or partially by another type of cation utilizing ion exchange techniques in a conventional manner. By means of such cation exchange, it has been possible to vary the properties of a given aluminosilicate by suitable selection of the cation. The spaces between the tetrahedra are occupied by molecules of water prior to dehydration.\nPrior art techniques have resulted in the formation of a great variety of synthetic aluminosilicates. These aluminosilicates have come to be designated by letter or other convenient symbols, as illustrated by zeolite A (U.S. Pat. No. 2,882-243), zeolite X (U.S. Pat. No. 2,882,244), zeolite Y (U.S. Pat. No. 3,130,007), zeolite ZK-5 (U.S. Pat. 3,247,195), zeolite ZK-4 (U.S. Pat. No. 3,314,752) and zeolite ZSM-5 (U.S. Pat. No. 3,702,886) merely to name a few.\nApplicants know of no prior art methods of zeolite preparation utilizing the present improvement. In fact, the present improved method of zeolite preparation is distinctly different from the current practice in the synthesis of zeolites having a silica/alumina mole ratio greater than 10."}
-{"text": "Diabetes refers to a disease process derived from multiple causative factors and characterized by elevated levels of plasma glucose or hyperglycemia in the fasting state or after administration of glucose during an oral glucose tolerance test. Persistent or uncontrolled hyperglycemia is associated with increased and premature morbidity and mortality. Often abnormal glucose homeostasis is associated both directly and indirectly with alterations of the lipid, lipoprotein and apolipoprotein metabolism and other metabolic and hemodynamic disease. Therefore patients with Type 2 diabetes mellitus are at especially increased risk of macrovascular and microvascular complications, including coronary heart disease, stroke, peripheral vascular disease, hypertension, nephropathy, neuropathy, and retinopathy. Therefore, therapeutical control of glucose homeostasis, lipid metabolism and hypertension are critically important in the clinical management and treatment of diabetes mellitus.\nThere are two generally recognized forms of diabetes. In Type 1 diabetes, or insulin-dependent diabetes mellitus (IDDM), patients produce little or no insulin, the hormone which regulates glucose utilization. In Type 2 diabetes, or noninsulin dependent diabetes mellitus (NIDDM), patients often have plasma insulin levels that are the same or even elevated compared to nondiabetic subjects; however, these patients have developed a resistance to the insulin stimulating effect on glucose and lipid metabolism in the main insulin-sensitive tissues, which are muscle, liver and adipose tissues, and the plasma insulin levels, while elevated, are insufficient to overcome the pronounced insulin resistance.\nInsulin resistance is not primarily due to a diminished number of insulin receptors but to a post-insulin receptor binding defect that is not yet understood. This resistance to insulin responsiveness results in insufficient insulin activation of glucose uptake, oxidation and storage in muscle and inadequate insulin repression of lipolysis in adipose tissue and of glucose production and secretion in the liver.\nThe available treatments for Type 2 diabetes, which have not changed substantially in many years, have recognized limitations. While physical exercise and reductions in dietary intake of calories will dramatically improve the diabetic condition, compliance with this treatment is very poor because of well-entrenched sedentary lifestyles and excess food consumption, especially of foods containing high amounts of saturated fat. Increasing the plasma level of insulin by administration of sulfonylureas (e.g. tolbutamide and glipizide) or meglitinide, which stimulate the pancreatic \u03b2 cells to secrete more insulin, and/or by injection of insulin when sulfonylureas or meglitinide become ineffective, can result in insulin concentrations high enough to stimulate the very insulin-resistant tissues. However, dangerously low levels of plasma glucose can result from administration of insulin or insulin secretagogues (sulfonylureas or meglitinide), and an increased level of insulin resistance due to the even higher plasma insulin levels can occur. The biguanides increase insulin sensitivity resulting in some correction of hyperglycemia. However, the two biguanides, phenformin and metformin, can induce lactic acidosis and nausea/diarrhea. Metformin has fewer side effects than phenformin and is often prescribed for the treatment of Type 2 diabetes.\nThe glitazones (i.e. 5-benzylthiazolidine-2,4-diones) are a more recently described class of compounds with potential for ameliorating many symptoms of Type 2 diabetes. These agents substantially increase insulin sensitivity in muscle, liver and adipose tissue in several animal models of Type 2 diabetes resulting in partial or complete correction of the elevated plasma levels of glucose without occurrence of hypoglycemia. The glitazones that are currently marketed are agonists of the peroxisome proliferator activated receptor (PPAR), primarily the PPAR-gamma subtype. PPAR-gamma agonism is generally believed to be responsible for the improved insulin sensititization that is observed with the glitazones. Newer PPAR agonists that are being tested for treatment of Type II diabetes are agonists of the alpha, gamma or delta subtype, or a combination of these, and in many cases are chemically different from the glitazones (i.e., they are not thiazolidinediones). Serious side effects (e.g. liver toxicity) have occurred with some of the glitazones, such as troglitazone.\nAdditional methods of treating the disease are still under investigation. New biochemical approaches that have been recently introduced or are still under development include treatment with alpha-glucosidase inhibitors (e.g. acarbose) and protein tyrosine phosphatase-1B (PTP-1B) inhibitors.\nCompounds that are inhibitors of the dipeptidyl peptidase-IV (\u201cDPP-4\u201d) enzyme are also under investigation as drugs that may be useful in the treatment of diabetes, and particularly Type 2 diabetes. See WO 97/40832; WO 98/19998; U.S. Pat. No. 5,939,560; U.S. Pat. No. 6,303,661; U.S. Pat. No. 6,699,871; U.S. Pat. No. 6,166,063; Bioorg. Med. Chem. Lett., 6: 1163-1166 (1996); Bioorg. Med. Chem. Lett., 6: 2745-2748 (1996); Ann E. Weber, J. Med. Chem., 47: 4135-4141 (2004); D. Kim, et al., J. Med. Chem., 48: 141-151 (2005); and K. Augustyns, Exp. Opin. Ther. Patents, 15: 1387-1407 (2005). The usefulness of DPP-4 inhibitors in the treatment of Type 2 diabetes is based on the fact that DPP-4 in vivo readily inactivates glucagon like peptide-1 (GLP-1) and gastric inhibitory peptide (GIP). GLP-1 and GIP are incretins and are produced when food is consumed. The incretins stimulate production of insulin. Inhibition of DPP-4 leads to decreased inactivation of the incretins, and this in turn results in increased effectiveness of the incretins in stimulating production of insulin by the pancreas. DPP-4 inhibition therefore results in an increased level of serum insulin. Advantageously, since the incretins are produced by the body only when food is consumed, DPP-4 inhibition is not expected to increase the level of insulin at inappropriate times, such as between meals, which can lead to excessively low blood sugar (hypoglycemia). Inhibition of DPP-4 is therefore expected to increase insulin without increasing the risk of hypoglycemia, which is a dangerous side effect associated with the use of insulin secretagogues.\nDPP-4 inhibitors also have other therapeutic utilities, as discussed herein. DPP-4 inhibitors have not been studied extensively to date, especially for utilities other than diabetes. New compounds are needed so that improved DPP-4 inhibitors can be found for the treatment of diabetes and potentially other diseases and conditions. In particular, there is a need for DPP-4 inhibitors that are selective over other members of the family of serine peptidases that includes quiescent cell proline dipeptidase (QPP), DPP8, and DPP9 (see G. Lankas, et al., \u201cDipeptidyl Peptidase-IV Inhibition for the Treatment of Type 2 Diabetes,\u201d Diabetes, 54: 2988-2994 (2005). The therapeutic potential of DPP-4 inhibitors for the treatment of Type 2 diabetes is discussed by D. J. Drucker in Exp. Opin. Invest. Drugs, 12: 87-100 (2003); by K. Augustyns, et al., in Exp. Opin. Ther. Patents, 13: 499-510 (2003); by J. J. Holst, Exp. Opin. Emerg. Drugs, 9: 155-166 (2004); by H.-U. Demuth in Biochim. Biophys. Acta, 1751: 33-44 (2005); by R. Mentlein, Exp. Opin. Invest. Drugs, 14: 57-64 (2005)"}
-{"text": "1. Field of the Invention\nThe present invention relates to a segmented brake light, and more particularly to a brake light composed of multiple light modules that are selectively and progressively activated according to different pressure a driver stepped on the brake pedal.\n2. Description of Related Art\nA top-installed rear brake light near the rear windshield in a vehicle is usually started when the driver is stepping on a brake pedal to warn other people behind the vehicle to avoid any collision.\nHow the top-installed rear brake light works is often determined by the sensitivity of the brake light and a driver's behaviors. A high sensitivity light can be easily activated with a small pressure when the driver slightly steps on the brake pedal. However, the driver may just move the feet to the brake pedal as a preparatory action for braking, not mean to brake the car. Comparing to the high sensitive operations, a low sensitively brake light requires a deep stepping or high pressure on the brake pedal. For other drivers behind the car, they may have no sufficient time in emergency brakes. The complete turn-off and turn-on operations of a conventional light are unable to provide other drivers with a good pre-warning effect.\nHowever, to produce a top-installed rear brake light that can display multiple effects according to different statuses, some complex factors such as modifying original circuits in the vehicle may discourage consumers from equipping the vehicle with such a brake light.\nTherefore, the present invention provides a new segmented brake light to overcome the foregoing drawbacks."}
-{"text": "1. Technical Field\nThe present disclosure relates to methods and valves for controlling the flow of fluid through a bore and more particularly, the disclosure relates in some embodiments to methods and ball valves for use in the oil and chemical process industry. More particularly, this disclosure relates to double piston, trunnion-mounted ball valves.\n2. Background Art\nBall valves are commonly used in both the oil and chemical process industries. A type of ball valve used to control flow of a fluid is an apertured ball valve such as is disclosed in PCT Patent Application No. WO 93/03255 published on Feb. 18, 1993, incorporated by reference herein. In an apertured ball valve the valve operation or function may be broken down into two separate stages. Firstly, the ball moves between an open and a closed position by rotating through 90 degrees such that the ball aperture from an orientation coaxial with the flow direction, i.e. when the valve is open, to a position whereby the ball aperture is normal or perpendicular to the flow direction and the valve is closed. Secondly, the valve seals in the closed position to prevent flow through the bore across the ball valve. Therefore, the on-off control of flow through the valve is achieved by rotating the ball through 90\u00b0 within the valve housing. Another ball valve is disclosed in U.S. Pat. No. 6,708,946, the teachings of which are incorporated herein by reference.\nThere are two basic types of ball valve mechanisms which currently exist. First, there is the trunnion mounted ball system in which the ball element is positionally constrained inside the valve, usually by radial bearings. The ball is rotated by the application of torque through a valve stem to the trunnion. Sealing occurs as a result of the valve seat on the upstream (or high pressure) side of the valve \u201cfloating\u201d onto the ball element and causing engagement between a surface of the valve seat and the surface of the ball. The advantage of this system is that it provides highly reliable rotation between the valve open and the closed positions. The principal disadvantage of this system is that seal reliability is reduced because the sealing force only develops in proportion to the annular area of the valve seat. Further, in high pressure applications, the force exerted on the ball on the upstream side of the valve can result in deformation of the ball and leakage between the ball and another valve seat located on the downstream (or lower pressure) side of the valve. Thus, when trunnion mounted ball systems are used in high pressure wells and especially those in which the well fluid has a high proportion of particulate matter, being generally known as \u201caggressive\u201d wells, the pressure is such that fluids and/or particulate matter may leak past seals between the ball and the valve seats. This often results in the valve not achieving integrity of sealing. In such cases, this type of ball valve is unable to operate properly in such conditions.\nThe second type of ball valve mechanism which effects the abovementioned function is known as the \u201cfloating ball system\u201d. In this system the ball is not positionally constrained relative to the valve body. Rotation is caused by the application of force to a point which is offset from the ball centre which, in conjunction with the mating curvatures of the ball and seat, cause the ball to rotate. Sealing occurs as a result of the ball \u201cfloating\u201d onto the valve seat. The advantage of this mechanism is that the reliability of the seal is increased, because the sealing force develops in proportion to the circular area of the ball to seat contact. The disadvantage of this type of mechanism is that the rotational reliability is reduced as the friction factor between the ball and seat are considerably larger than that of trunnion mounted devices. With high pressure and aggressive types of wells and particulate flows of the type described above, the reliability of this valve in those applications creates a problem in that the torque necessary to rotate the ball becomes excessively high, and thus, the valve can seize between the open and the closed position giving rise to serious problems in both operational and safety terms.\nIt would be desirable to provide a method and/or improved ball valve design which may obviate or mitigate at least one or more of the aspects associated with the aforementioned disadvantages."}
-{"text": "Immunoglobulin compositions prepared from human plasma and suitable for intravenous administration are known in the art and for several decades have played an important role in the treatment of a wide range of diseases. Immunoglobulins are used, for example, for the treatment of infections in humans and can be assigned to various classes with various biochemical and physiological properties. Immunoglobulin G participates in defending against viral antigens, whereas IgM is predominantly active in antibacterial and antitoxin immune responses.\nThe immunoglobulin solutions comprise IgG, IgA and IgM in various percentages, with different preparations having different treatment applications, e.g. preparations with a higher percentage of IgM are used in the prophylaxis or treatment of bacterial infections.\nThe immunoglobulin solutions are usually prepared from fractions of blood plasma or serum, e.g. Cohn fractions. These fractions are then subjected to a number of purification steps to remove contaminants such as viruses, denatured proteins, proteases and lipids.\nHuman plasma for fractionation is collected from thousands of donors and may contain pathogen viruses despite thorough testing of the source plasma. Therefore process steps to inactivate or remove viruses are essential in order to achieve safe products for use in medicine. Several techniques for virus inactivation/removal are known in the art, e.g. chemical treatments, irradiation with UVC light or nanometer filtration, which are performed in order to ensure overall virus safety.\nThe virus removal or inactivation capacity of the process steps is validated using laboratory scale models of the production process and for each step a removal or inactivation factor is determined. An increase of the inactivation/removal factor adds additional viral safety to the pharmaceutical product. Today guidelines from regulatory authorities require at least two effective steps for enveloped and non-enveloped viruses in the manufacture of plasma-derived pharmaceuticals. Although several methods, such as solvent/detergent treatment, octanoic acid treatment, nanometer filtration and heat treatment, are effective to inactivate or remove enveloped viruses there are only a few methods known to inactivate or remove non-enveloped viruses, for example Parvo viruses. These non-enveloped viruses are mostly very small, usually passing through nanometer filters with pore sizes above 20 nm. This pore size is too small for IgM molecules having a diameter up to 30 nm. Non enveloped viruses are effectively inactivated by chemicals like \u03b2-propiolactone which, however, also leads to a modified immunoglobulin with impaired functions. Another effective treatment is UVC-irradiation (EP1842561, CAF-DCF). However, known solvent/detergent treatments, octanoic acid treatment and mild heat treatment have no substantial effect on non-enveloped viruses.\nAs mentioned above, in addition to viruses which are potentially present it is also necessary to remove other contaminants like lipids, proteases, protein aggregates, and denatured immunoglobulins. The removal of all these contaminants is essential (1) to ensure the product complies with bio-safety guidelines regarding viral contamination, (2) in order for the product to be tolerated by the patient after intravenous administration, (3) to allow the product to be stable during long-term storage (any residual proteolytic activity might lead to degradation of the product over long-term storage, e.g. 2 years), and (4) to generate the desired compound mixture/pharmaceutical composition.\nAt the same time, however, it is essential that the purification steps to remove the contaminants do not interfere with the immunoglobulin molecules, so that as far as possible these retain their normal biological activity and are retained at high yield in solution. This balance is difficult to achieve since many known purification steps can also have a negative impact on the activity of the immunoglobulins, and in particular on IgM; for example extended irradiation times with UVC can reduce the yield of native and active IgM obtained in the final immunoglobulin solution. Not only does this lead to a reduction in efficacy of the final immunoglobulin solution but it can also cause the solution to be less well tolerated in vivo.\nAggregates and denatured immunoglobulins, the amount of which can be increased by certain purification steps, especially are a potential risk for the patients because they have a high capacity to activate complement unspecifically, leading to severe side effects in patients receiving these denatured immunoglobulins. Unspecific complement activation refers to the initiation of the complement cascade in the absence of specific antibody-antigen complexes. Unspecific complement activation is strictly to be avoided since it may cause undesirable side effects such as hypotension, flushing, headache, fever, chills, nausea, vomiting, muscle pain, dyspnoea and tachycardia. Specific complement activation, on the other hand, is desirable and it occurs only after the immunoglobulins have bound to their specific antigens.\nUnspecific complement activation is measured as the so called anticomplementary activity (ACA) by a standardized test described in the European Pharmacopoeia.\nThe role of the complement system in the immune defense of pathogens is well known. The complement system consists of about 20 proteins, which are activated sequentially. The classical complement pathway typically requires a specific antigen antibody complex for activation, whereas the alternative pathway can be activated by antigens without the presence of antibodies. The classical and the alternative pathway of complement activation all generate a protease C3-convertase. The C3-convertase cleaves and activates component C3, creating C3a and C3b, and causing a cascade of further cleavage and activation events over C5 convertase to C5a and C5b. C5b initiates the membrane attack pathway, which results in the membrane attack complex, consisting of C5b, C6, C7, C8, and polymeric C9. This is the cytolytic endproduct of the complement cascade which forms a transmembrane channel, which causes osmotic lysis of the target cells like bacteria.\nComplement activation additionally results in the formation of anaphylatoxins, including the biologically active protein C5a. This anaphylatoxin is a potent chemotactic agent for immune and inflammatory cells and induces cell activation and causing the release of histamine from mast cells. In situations of excessive or uncontrolled and/or unspecific complement activation, the overproduction of C5a can cause deleterious effects to patients.\nC5a is an effective leukocyte chemoattractant, causing the accumulation of white blood cells, especially neutrophil granulocytes, at sites of complement activation. C5a activates white blood cells and is a powerful inflammatory mediator. Whereas these functions are beneficial during specific antibody-antigen complex reactions all unspecific generation of C5a has to be avoided due to the potential side effects.\nUnspecific complement activation is a particular issue for IgM immunoglobulin preparations (i.e. those comprising at least 5% IgM) as in contrast to IgG preparations IgM antibodies easily aggregate in solution. IgM preparations are difficult to stabilize especially if they are enriched compared to plasma concentrations and stored in liquid solution. It is also known that IgM is a vigorous activator of complement; a single molecule bound to an antigen can activate complement. This is in contrast to IgG, where two or more molecules of IgG must be bound to an antigen in close association with each other to activate complement.\nStill further, the main indications treated by IgM containing immunoglobulin preparations are bacterial infections and sepsis. As these patients are already suffering from hypotension an additional unwanted generation of unspecific complement activation and C5a would lead to a clinical worsening of the patient's condition. Accordingly, IgM preparations have been described as being difficult to prepare for intravenous application.\nThere are several methods described in the art for the production of IgM containing immunoglobulin preparations from human plasma.\nThe initial purification of human IgM solutions has been carried out by classical Cohn plasma fractionation methods or its well known modifications (e.g. Cohn/Oncley, Kistler/Nitschmann). Using cold ethanol precipitation processes the IgM fraction is recovered in fraction III or fraction I/III (also called B or B+I). Starting from fraction III or I/III methods have been described for purification of protein solutions enriched in IgM. EP0013901 describes a purification method starting from fraction III including steps using octanoic acid, \u03b2-Propiolactone treatment and an adsorption step using an anionic exchange resin. This method is used to produce Pentaglobin\u00ae\u2014to date the only commercially available intravenous IgM product. \u03b2-propiolactone is a well known chemical used in sterilization steps in order to inactivate viruses which are potentially present. As \u03b2-propiolactone is a very reactive substance which causes the chemical modification of proteins there is also substantial loss of the anti-viral and anti-bacterial activities of the immunoglobulins. On the other hand this chemical modification results in an reduced anticomplementary activity compared to an chemically unmodified immunoglobulin. EP0352500 describes the preparation of an IgM concentrate for intravenous application with a reduced anti-complementary activity by using anionic exchange chromatography, \u03b2-Propiolactone, UVC light irradiation and an incubation step at increased temperature (40\u00b0 C. to 60\u00b0 C.). The preparation produced by this method was stable in liquid solution for a limited time due to the chemical modification. The IgM concentration was above 50% from the total immunoglobulin content.\nThe preparation of protein solutions enriched in IgM without chemical modification by \u03b2-propiolactone has been described in EP0413187 (Biotest) and EP0413188 (Biotest). These methods involve subjecting a suitable protein solution to octanoic acid treatment and anionic exchange chromatography, starting from Cohn fraction III or II/III. In patent EP0413187 (Biotest) the octanoic acid treatment is carried out by stirring for 15 min, in order to remove lipids being present in Cohn fraction III.\nThe preparation according to EP0413187 had a low anticomplementary activity, between 0.6 and 0.8 CH50/mg protein, but had to be stabilized and virus inactivated by \u03b2-propiolactone. Low anticomplementary activity is considered to be \u22661 CH50/mg protein according to EP monograph for immunoglobulins.\nEP0413188B1 (Biotest) describes the preparation of an IgM-enriched preparation for intravenous administration by using an anion exchange chromatography in order to reduce the anti-complementary activity. Additionally a heat treatment at pH 4-4.5 at 40 to 60\u00b0 C., preferably between 50 and 54\u00b0 C., was described to reduce the anticomplementary activity. This preparation had to be lyophilized to ensure stability of the preparation for several months. Long term stability as a liquid solution could not be shown.\nM. Wickerhauser et al. \u201cLarge Scale Preparation of Macroglobulin\u201d, Vox Sang 23, 119-125 (1972) showed that IgM preparations isolated by PEG precipitation had high anticomplementary activity (ACA) by a standard complement fixation test and this ACA activity was reduced 10 fold by incubating the IgM preparation at pH 4.0 at 37\u00b0 C. for 8 hours followed by readjustment to neutral pH. It was not demonstrated if this 10 fold reduction is sufficient to ensure intravenous tolerability. The authors did not assess the specific complement activating potential of their IgM concentrate, nor did they assess safety in any animal or human model.\nAnother method describes the use of mild-heat treatment of IgM preparations at 40 to 62\u00b0 C., preferably 45 to 55\u00b0 C., at pH 4.0 to 5.0 (EP 0450412, Miles) to reduce the unspecific complement activation. In this patent application octanoic acid is added to a Cohn fraction III suspension in order to remove prekallikrein activator and lipoproteins by centrifugation. Nevertheless this mild heat treatment led to partial loss of antigenic determinants of IgM. This may increase the risk of generating neo-antigens leading to a increased immunogenicity in humans or the loss of activity.\nThe preparation of an IgM containing protein solution for intravenous application by using a protease treatment (e.g. with pepsin) after an octanoic acid precipitation step has been described in EP0835880 (U.S. Pat. No. 6,136,312, ZLB). Protease treatment leads to partial fragmentation of the immunoglobulin molecule impairing the full functional activity of the Fab and Fc parts. Therefore protease-treated immunoglobulins cannot be regarded as unmodified. Also this preparation method leads to about 5% fragments with a molecular weight of <100 kD.\nThe described methods to carry out the octanoic acid treatment (EP0413187 and EP0835880) have the drawback that the octanoic acid treatment is not effective with respect to removal and inactivation of non-enveloped viruses, and does not remove substantially all proteolytic activity.\nIn EP 0345543 (Bayer, Miles) a highly concentrated IgM preparation with at least 33% IgM for therapeutic use is disclosed, the preparation being substantially free of isoagglutinin titres. In this patent application an octanoic acid precipitation is carried out by adding the octanoic acid and the isoagglutinins are removed by Synsorb affinity chromatography. The final preparation had to be freeze dried.\nAltogether the production of an IgM containing preparation with low anticomplementary activity is possible if the immunoglobulins are chemically or enzymatically modified and/or further purified by chromatography and/or subjected to a mild heat treatment. However, these methods have their drawbacks in the lack of virus removal/virus inactivation (and therefore virus safety), reduction in the amount of immunoglobulin molecules in native form and/or residual anticomplementary activity. As such, there is still a need to provide improved IgM containing immunoglobulin preparations suitable for intravenous administration in humans."}
-{"text": "Remote memory, prepaid accounts for use in purchasing goods and services are generally well known. Presently known schemes typically include a printed transaction card, for example a wallet-sized plastic or cardboard card, which bears on one side a unique authorization or account number and instructions for access to funds, services, and the like. Such prepaid cards have been used extensively throughout the world. One such example is the use of these remote memory cards as prepaid long distance telephone calling cards.\nIn contrast to stored value cards (e.g., \"smart cards\") wherein a remaining account balance is stored within a memory resident in the card, remote memory systems generally store information pertaining to a prepaid account at a central host computer. The host computer typically stores information relating to the available balance remaining in the account, as well as information pertaining to past activity associated with the account. In particular, the host computer may store transaction data relating to various goods or services purchased using the card. In the context of a prepaid telephone calling card, the host computer may store telephone call data, including the date, time, duration, and various other parameters relating to calls which were placed using the prepaid telephone card. The host computer may be accessed via a telephone or data line by the consumer through the use of an authorization code, personal identification number (PIN), or the like.\nThe use of prepaid remote memory telephone cards is particularly prevalent in the United States. A typical prepaid telephone calling card includes a toll-free telephone number used by a consumer to access a host computer system, a unique authorization code associated with the card and a corresponding remote account, and dialing instructions. When a customer desires to use the card to place a long distance call, he dials the toll-free number, thereby accessing the host system which manages the remote accounts. The consumer then enters a predetermined authorization number for allowing access to the database. Next, the consumer enters the desired long distance telephone number, and the system connects the consumer with the desired calling destination. Long distance telephone charges attributed to the telephone call are deducted from the remaining balance in the account and the call is terminated when the account is fully consumed. A call history, which includes information pertaining to the calls charged to the various accounts, may also be maintained by the host computer for each account.\nIn the United States, preprinted instructions on how to use the transaction cards are generally printed directly on the cards in a single language, such as English. The use of such cards by individuals who are not conversant in the English language can be inconvenient or even impossible, depending on their fluency in the English language and their familiarity with how to use the transaction cards. While instructions may be printed on a single card in several languages, the number of languages that may be printed on a card is directly related to the size of the card and the level of detail to which the instructions will be printed; in the event wallet-sized cards are used, limited space usually precludes more than one language being printed on the card. Another option in attempting to resolve this language problem is to print a series of cards in a variety of different languages. However, this option has the disadvantage, among others, of increasing the issuer's and merchant's inventory by a factor equal to the number of languages printed. A third option is to assign each language a separate telephone number. Once connected to the host computer by phone, the consumer can be instructed on how to use the transaction card in that particular language. This however, has the disadvantage of requiring the issuer to maintain a phone number and connection for each language.\nConsequently, a system and method is thus needed which will overcome the shortcomings of the prior art and will allow the issuer to easily produce transaction cards that may be used by consumers in a plurality of different languages."}
-{"text": "(1) Field of the Invention\nThe present invention relates to encapsulated foods filled with a desired edible liquid and the process for preparing drinks which include such encapsulated foods. Furthermore, this invention relates to a process for preparing encapsulated foods to make, encapsulated dressings, encapsulated alcohol drinks, encapsulated coffee or red tea etc. by using dressings, fruit juice, alcoholic drinks and other liquids, popular foods, liquid sweetening materials, etc. as edible liquids and the process for preparing drinks which include these encapsulated foods.\nGenerally, popular foods such as alcohol, coffee, red tea, cocoa, etc. as well as fruit juice are used as drinks because of their liquid form. Moreover, acidic condiments, i.e. dressings and condiment vinegar etc., being normally liquid, are poured over a salad, for example, by using a spoon. However, the poured condiments either permeate into the salad itself or are dispersed and absorbed into any food surrounding it. Therefore, it is impossible to visually confirm the quantity that has been used or to know for certain how much of it enhanced the flavor of the salad. Although powdered condiments and granular ones may resolve these problem to a certain extent, the use of powdered ones is somewhat restricted due to their quality or taste.\nThe present invention intends to resolve these problem. More specifically, it provides a new process for preparing edible dressings, edible alcoholic drinks, edible coffee or tea, etc. which have the taste and quality of liquid dressings, fruit juices, alcoholic drinks, popular liquid foods, liquid sweetening materials, etc., and besides have more visual character and are not too soft.\n(2) Description of Prior Preparation\nPreviously, a process for preparing capsules using solutions of alginic acid salt and calcium salt is conducted by dropping solutions of alginic acid salt and calcium into another solution to form a membrane of calcium alginate, thereby forming the capsule. However, it is impossible to obtain suitable capsules when this process is applied as it is to dressings, alcoholic drinks or to any preferred drinks.\nFor example, when an acidic condiment per se is subjected to encapsulation, the equilibrium of the alginic acid ion and sodium ion is lost by the acidifying the sodium alginate solution to agglomerate and generate alginic acid separately. Also, the calcium ions for forming the membrane are fewer because they react with acid so as to extract calcium salt by acidifying the calcium salt solution; therefore, in any way, it is difficult to produce good capsules.\nIt may be considered that the process is to first form the capsules and then fill them with an acidic condiment liquid. However, in this case, the capsules consisting of calcium alginate are influenced by the acid due to a time lapse, thus causing various problems to arise such as how to produce capsules with strong membranes which will neither collapse nor change in quality under preservation when filling them with the acidic condiment liquid and how to fill the acidic condiment liquid into the inner spaces of the capsules already formed.\nWhen fruit juice per se is encapsulated as a core material, various problems arise such as that fruit juice reacts with alginic acid to increase its viscosity and alginic acid generates and agglomerates separately since fruit juice is normally an organic acid liquid of pH 3-4. Furthermore, it becomes difficult to conduct a strong bridge reaction in view of the competing COOH group of alginic acid molecules and the COOH group of fruit juice molecules because fruit juice reacts with calcium salt to decrease the number of calcium ions and further generates alginic acid calcium, producing an inferior taste so that, being not edible, this causes a problem that is difficult to solve.\nConcerning alcoholic drinks to be encapsulated as core material, here also various problems arise making it difficult to form the desired capsules despite the addition of products such as xanthene gum, etc.; besides, drops of alcoholic drinks collapse on the surface of alginic acid salt liquid when they are dropped into it. The reason for such a phenomena is due to the fact that alcoholic drinks per se are low in molecular compounds and have a low viscosity. In any case, there is also the problem of solving the difficulties in forming capsules by such a reaction.\nThe same problem exists in the case of other liquids and liquid sweeters. Namely, when these substances are encapsulated as core material those in which the mixture of alginic acid salt with liquids or liquid sweetening material is dropped into calcium salt water solution may be encapsulated well; however, this manner of preparation is not desirable for taste because the resulting capsules become hard. Furthermore, when the mixture of calcium salt water solution and liquid sweetening material is dropped into an alginic acid water solution, an unsuitable taste is generated by the remaining calcium constituent in the sweet materials. In any case, it is difficult to prepare suitable capsules. Such problems also occur with encapsulated foods including cut solid foods containing the core liquid."}
-{"text": "The subject matter discussed in the background section should not be assumed to be prior art merely as a result of its mention in the background section. Similarly, a problem mentioned in the background section or associated with the subject matter of the background section should not be assumed to have been previously recognized in the prior art. The subject matter in the background section merely represents different approaches, which in and of themselves may also correspond to implementations of the claimed technology.\nThe Naive Bayes Classifier (NBC) is a stalwart of the machine-learning community, often the first algorithm tried in a new arena, but also one with the inherent weakness of its hallmark assumption\u2014that features are independent. Bayesian networks relax this assumption by encoding feature dependence in the structure of the network. This can work well in classification applications in which there is substantial dependence among certain sets of features, and such dependence is either known or is learnable from a sufficiently large training set. (In the latter case, one may use any of a number of structure-learning algorithms; they don't always work well because the structure-learning problem is very difficult.) Undirected graphical models are a different way to relax this assumption. Like Bayesian networks, they also require domain knowledge or a structure-learning procedure. Domain knowledge can be hand coded into the structure, but this only works when domain knowledge is available and even then hand coding is usually very laborious. Alternatively, structure-learning algorithms require extensive training data and present a computationally complex problem, in addition to the concern of over-fitting a model using a limited data set, thereby creating a model that predicts structures that are less likely to work well on data not seen during training Undirected graphical models may also fail to work well when either the domain knowledge is wrong or when the instance of the structure-learning problem is either intrinsically hard or there is insufficient data to train it.\nThe technology disclosed relates to machine learning (ML) systems and methods for determining feature dependencies\u2014methods which occupy a space in between NBC and Bayesian networks, while maintaining the basic framework of Bayesian Classification.\nBig data systems now analyze large data sets in interesting ways. However, many times systems that implement big data approaches are heavily dependent on the expertise of the engineer who has considered the data set and its expected structure. The larger the number of features of a data set, sometimes called fields or attributes of a record, the more possibilities there are for analyzing combinations of features and feature values.\nAccordingly, an opportunity arises to automatically analyze large data sets quickly and effectively. Examples of uses for the disclosed classification systems and methods include identifying fraudulent registrations, identifying purchase likelihood for contacts, and identifying feature dependencies to enhance an existing NBC application. The disclosed classification technology consistently and significantly outperforms NBC."}
-{"text": "Direct conversion receivers (DCRs) are prevalent in mobile communication systems due to their simplicity, low cost and for the reason that much of the signal processing can be achieved in the digital domain. A drawback to the elimination of an intermediate frequency stage is an increased susceptibility to second-order intermodulation effects. For example, intermodulation products generated from strong interfering signals can be downconverted to baseband with relatively weaker signals of interest, thereby desensitizing the DCR to the desired signal. Such interfering signals, referred to herein as blocking signals, are ubiquitous; they may originate from communication signals in adjacent channels and/or from transmission sources that, even when far removed in frequency from the channel of interest, transmit at a power level sufficient to cause significant interference problems in a DCR, including bit error rate (BER) and/or signal-to-noise ratio (SNR) degradation.\nIn view of the susceptibility to distortion by second-order effects, design specification of DCRs typically include high input-related second order intercept point (IIP2) criteria. IIP2 is the theoretical input signal level at which second-order intermodulation products are equal in power to that of a desired signal. Thus, if the IIP2 can be made higher, the power of an interfering signal must reach a correspondingly higher level to have an equivalent detrimental effect on the DCR. It is thus clearly desirable to establish IIP2 at as high a level as possible.\nSecond-order intermodulation distortion is an effect of non-linear behavior of components in a DCR and downconverting mixers are often mostly responsible for limits on the level to which IIP2 can be established. Compensation techniques can be applied to increase the IIP2, but changes in temperature and/or frequency can reduce the IIP2 to effectively undo applied compensation measures. Indeed, an IIP2 established at 100 dBm can be reduced to 50 dBm in response to a change in frequency or temperature.\nGiven the potential improvements in DCR performance that can be expected upon overcoming the foregoing limitations in the art, substantial resources are increasingly devoted toward developing interference-tolerant designs in direct conversion receivers."}
-{"text": "1. Field of the Invention\nThe present invention relates to a head protecting airbag system that deploys an airbag downward along the side of a vehicle body in a passenger compartment when the vehicle is involved in either a side collision or a rollover. In particular, the present invention relates to a head protecting airbag system that deploys an inflation section between an occupant's head and the vehicle body.\n2. Description of the Related Art\nIn the recent trends, a head protecting airbag system is mounted in vehicles as a supplemental restraint system. The head protection airbag system has a curtain airbag that is deployed downward from a roof side rail when the vehicle is involved in a side collision or rollover. As one of such head protecting airbag systems, there has been proposed a head protecting airbag system designed to deploy the airbag between the head of the occupant and the side of the vehicle body in order to protect the occupant's head.\nThis type of head protecting airbag system is described in Japanese Patent Application Publication No. 2004-58848 (JP-A-2004-58848). In this head protecting airbag system, the airbag includes a gas supply path and a main inflation part disposed on the lower side of the gas supply path. The airbag is folded with the main inflation part rolled up toward the outer side in the vehicle and with the gas supply path not rolled, but folded for the sake of easier deployment of the airbag upon gas supply. In addition, in Japanese Patent Publication No. 3656156, at the initial stage of deployment of the airbag, a part of the airbag deploys into a thin plate shape between the pillar and the headrest of the occupant seat with no gas supplied. Then, gas is supplied into the thin plate airbag from an inflow port formed on the lower side of the airbag. This prevents the inflating airbag from being interfered with by the headrest of the occupant seat.\nHowever, in order to immediately deploy the airbag in the longitudinal direction of the vehicle along the side of the passenger compartment, the airbag has to initially deploy inward towards the passenger compartment over the top end of the pillar garnish."}
-{"text": "1. Field of the Invention\nThe present invention relates to a medical image display control device, a method of operation for the same, and a medical image display control program for displaying the internal tissue of a subject so as to be superimposed on a three-dimensional image of the subject.\n2. Description of the Related Art\nIn the related art, the planning of actual surgery or the like has been performed using a three-dimensional image captured before surgery. When performing such surgical planning, the internal tissue, such as blood vessels at the back of the skin tissue or the bone tissue, needs to be visualized in some cases.\nSpecifically, for example, in the case of performing craniotomy, only things on the surface are visible when performing actual surgery since the head has a layer structure of a skin region, a bone region, and a brain region from the surface. Therefore, when incising the skin, blood vessels under the skin are the important tissue for the surgery. In addition, when opening the bone, brain veins or brain arteries under the bone are the important tissue for the surgery.\nIn such a case, in general, the internal tissue, such as blood vessels, can be displayed by extracting the skin tissue or the bone tissue and adjusting the opacity or the attenuation.\nIn order to display the internal tissue, JP2007-325629A has proposed to display a three-dimensional image of the internal tissue excluding, for example, a region of the skin tissue by setting a predetermined range from the surface of the subject as a removal layer and constructing a three-dimensional image from which the removal layer has been removed."}
-{"text": "1. Field of the Invention\nThis invention relates generally to shaft seals that are used in hydrogen inner-cooled turbine-generators, and more particularly to methods and apparatus for cooling the shaft seals in such turbine-generators.\n2. Statement of the Prior Art\nConductor cooling is a conventional process used in very large turbine-generator systems for dissipating the armature and field coil losses of the turbine-generators to cooling media within their coil insulation wall. The turbine-generators using such conductor cooling are also referred to variously as \"inner-cooled,\" \"supercharged,\" or \"direct-cooled\", and the cooling medium that is most often used in such turbine-generators is hydrogen.\nHydrogen inner-cooled turbine-generators typically operate in a pressurized hydrogen atmosphere that provides cooling for all of the turbine-generator except, in some instances, its armature coils. Such hydrogen inner-cooled turbine-generators are often operated at 60 lbf/in.sup.2 (i.e., about 4219 grams per square centimeter) or more in order to increase the mass flow of the hydrogen, and to reduce its temperature rise.\nIn order to minimize leakage of pressurized hydrogen cooling medium from an operational turbine-generator, shaft seals are typically used in hydrogen inner-cooled turbine-generators for maintaining an oil film under pressure in a small clearance between the rotating shaft of the turbine-generator and a stationary member surrounding the shaft at both ends of the turbine-generator. The construction of such shaft seals may be similar to a journal bearing with a cylindrical oil film or similar to a spring-loaded thrust bearing with the oil film in a plane at right angles to the shaft axis. In either case, the oil film is maintained by an oil supply pressure that is higher than the hydrogen pressure.\nOils used in such shaft seals can absorb about 10% by volume of either hydrogen or air. It is important that the flow of oil in those shaft seals toward their hydrogen side be minimized in order to reduce both the amount of air that is carried into the hydrogen inner-cooled turbine-generator and the amount of hydrogen that is carried out. Moreover, it is important to minimize temperature differences between the \"hydrogen side\" of the shaft seal and its \"air side\" so that differential thermal expansion of the seal ring which comprises the stationary member can be minimized.\nPrior art shaft seals have typically employed separate supplies of oil for their hydrogen side and their air side, each such oil supply including a heat exchanger operated in one of two general fashions. One method selects a cooling media (e.g., water), and carefully sizes the heat exchanger of each oil supply to provide for open loop control of seal oil temperatures at the outlets on either side of the shaft seals. Any problems with heat exchanger fouling, decreased oil flow or decreased water flow when experienced with this method requires recognition and subsequent correction by an operator.\nAnother method utilizing separate heat exchangers for the hydrogen side and the air side of prior art shaft seals provides automated control valves to sense the temperature difference between the hydrogen side oil supply and the air side oil supply, and operates the particular heat exchanger in the oil supply having the higher temperature. Not only is this other method complicated in its design by virtue of the additional components that are required (e.g., two heat exchangers and many automated control valves), but it is also difficult in its implementation due to the necessity for precise, reliable components to accomplish such control over a very narrow temperature range."}
-{"text": "1. Field of the Invention\nThe present invention relates to a fuel conditioner adapted to be placed in-line in a fuel delivery system and is adapted to condition fuel for improved purity, extended storage life, and reduced-engine wear.\n2. Description of the Related Art\nInternal combustion engine systems typically are provided fuel from a remote storage tank via fuel lines and the fuel is driven either by gravity or an active pump. The systems often include an in-line filter to remove particulate impurities. However, filters typically are passive devices that can only screen out particles above a certain size.\nFilters are relatively ineffectual against biological processes that often occur in fuel. As an example, diesel fuel often accumulates water in storage. The fuel, especially with water present, can support the growth of certain bacteria, fungi, and algae. It is known to add anti-biological agents such as bactericides and fungicides to the fuel, however this requires the undesirable additional effort of adding the anti-biological agents to the fuel.\nAn additional contaminant that can be present in fuel that is not particularly well handled by conventional filters is metal contamination. Small particles of metal can become entrained in the fuel from wear in fuel pumps and corrosion in fuel delivery systems. These small metallic particles can be too small to be effectively trapped by a filter, yet large enough to cause undesirable wear and deposits in the engine.\nFrom the foregoing, it can be appreciated that there is a need for a fuel conditioning system that can inhibit the growth of biological contaminants in fuel yet avoids the inconvenience of mixing additives with the fuel. There is also a need for a system to remove metallic particles from a fuel supply.\nThe aforementioned needs are satisfied by the present invention which, in one aspect, is an in-line fuel conditioner receiving a flow of liquid fuel the conditioner comprising a housing, end caps attached to ends of the housing and a plurality of reactive elements contained within the housing such that the fuel passes over the reactive elements and wherein the reactive elements comprise separate stainless steel, zinc, and copper members such that the overall composition of the reactive elements is approximately 50-40% stainless steel, 40-30% zinc, and 30-20% copper by weight and wherein the in-line fuel conditioner inhibits the growth of biological agents entrained within the fuel. In certain embodiments, the invention also includes a magnet wherein the magnet retains ferrous metal particles entrained within the fuel flow and in one embodiment, the magnet is positioned with the housing. In certain embodiments, the reactive elements are approximately 0.125xe2x80x3 in major dimension. In another aspect, the invention is an internal combustion engine system utilizing fuel and including an in-line fuel conditioner wherein the in-line fuel conditioner contains a plurality of reactive elements comprising separate stainless steel, zinc, and copper members wherein the reactive elements have a weight composition of 50-40% stainless steel, 40-30% zinc, and 30-20% copper and a magnet wherein the in-line fuel conditioner inhibits the growth of biological contaminants and retains ferrous particulates entrained within the fuel.\nThe in-line fuel conditioner of the present invention can be readily installed in an existing fuel delivery system using commonly available tools and known mechanical techniques. The fuel conditioner inhibits the growth of biological contaminants without the inconvenience of treating the fuel with additives. The fuel conditioner also retains ferrous particulates entrained within the fuel thereby reducing wear to an engine system so equipped. These and other objects and advantages will become more fully apparent from the following description taken in conjunction with the accompanying drawings."}
-{"text": "1. Field of the Invention\nThe present invention relates to an automatic document conveying device that conveys a document to a reading position automatically, and an image forming apparatus including the automatic document conveying device.\n2. Description of Related Art\nAs an automatic document conveying device used in a copier, a scanner, a facsimile, etc., there are well-known automatic document conveying devices in a form capable of inverting front and back surfaces of a document so as to read both surfaces of the document. In recent years, along with the downsizing of an image forming apparatus, there is also a demand for the downsizing of the automatic document conveying device as well as an image forming portion of the image forming apparatus.\nFIG. 13 illustrates a related technology of such an automatic document conveying device. The related technology of the automatic document conveying device includes a document stacking tray (101), a pickup roller (102), a sheet feeding roller (103), a conveying drum (104), a front reading roller (105) and a back reading roller (106) disposed so as to be in contact with the conveying drum (104), a branching member (107), a lower inversion roller (108), an upper discharge roller (109), an inversion guide (110), and a discharge tray (111) in this order in a conveying direction (direction of an arrow \u201ca\u201d) of a document (D). Further, a document reading position (R) is located below the conveying drum (104) between the front reading roller (105) and the back reading roller (106). For convenience of the description, an upper surface of the document placed on the document stacking tray (101) is referred to as an A-surface and a lower surface thereof is referred to as a B-surface.\nWhen only the A-surface of the document is read, the A-surface of the document (D) supplied by the pickup roller (102) from the document stacking tray (101) is read at the document reading position (R). The branching member (107) is supported so as to be swingable about an axis at its back end. When only the A-surface of the document (D) is read, the branching member (107) is placed at a position indicated by a solid line illustrated in the figure, and the document (D) is discharged to the discharge tray (111).\nWhen both the surfaces (A-surface and B-surface) of the document are read, after the A-surface of the document (D) conveyed from the document stacking tray (101) is read at the document reading position (R), the document (D) is inverted by the lower inversion roller (108), the upper discharge roller (109), and the inversion guide (110) due to the presence of the branching member (107) at a position indicated by a dotted line illustrated in the figure. Then, the document (D) is conveyed to the document reading position (R) and the B-surface is read. The document (D) is further inverted again due to the similar operation and discharged to the discharge tray (111).\nIn the above-mentioned related technology, the number of members such as a conveying roller for use in the automatic document conveying device is reduced by configuring the automatic document conveying device as described above. However, drive sources of the branching member (107) and the lower inversion roller (108) are provided separately, and hence there is a problem in that spaces are required for setting the respective drive sources. Further, in a case of using an inexpensive solenoid as a drive source, the noise from the drive source increases, which requires countermeasures against the noise. Consequently, there is also a problem in that the number of parts increases to increase the number of assembly steps, leading to an increase in production cost."}
-{"text": "This invention relates to a high-resolution, cryogenic, side-entry type specimen stage for a high-voltage electron microscope in which in-situ four point D.C. electrical resistivity measurements and electron optical observation can be made.\nStudies of the effect of irradiation on metals are of great importance today. For example, extensive work is being done on electron damage studies of metals which might be useful in fusion reactors. For this purpose optical studies using electron microscopes are of importance. Likewise, resistivity studies are of great importance because defect damage increases the electrical resistivity which can be detected before the defects become pronounced enough to see visually. While both types of studies have been made in the past, no instruments are available within which electron optical studies and electrical resistivity measurements can be made simultaneously at cryogenic temperatures. The electrical resistivity measurements would complement the electron optical observation in the high-voltage electron microscope to yield a unique opportunity to investigate defect production in metals by electron irradiation over a wide range of defect concentration. The specimen stage described herein is designed to achieve the same electron optical resolution (.about.5 A) below 10 K as can be obtained at room temperature."}
-{"text": "This invention relates in general to folding cartons and, in particular, to a folding carton which is adapted to be used for containing and packaging smaller individually pre-packaged items within the container.\nMore specifically, but without restriction to the particular uses which are shown and described, this invention relates to a carton formed from an open-ended tube in which the closures thereof may be formed from an integral portion of the tube body. Portions of the packaged items contained within the carton may be observed through the carton to determine such things, for example, as size, color or \"UPC\" coding.\nIn the retail trade, it is quite common to merchandise such articles as pre-packaged hoisery, smoking tobacco, office supplies, dry goods, drug supplies and other sundry items in individual packages to be purchased by a retail customer. Such items are frequently purchased by a retailer from a wholesale distributor in smaller quantities than in large case lots containing a greater quantity of these items. Therefore, it has become desirable to provide an intermediate packing container or carton within such larger packing cases or shipping containers. Such intermediate packing containers or cartons are referred to in the trade as a \"distributor pack\" and contain a smaller quantity of the individually pre-packaged consumer or retail items.\nDistributor packs used in this manner are especially useful for packing irregular shaped individual items. Such items, for example, as blister-packed tape dispensers are difficult to package as individual items in case lot quantities due to their irregular shape. Therefore, smaller quantites of such items are frequently packaged in a container, such as a distributor pack, which is used as an intermediate packer in a larger case to facilitate orderly packing and shipping of such items in case lot quantites. Such distributor pack cartons may also be utilized to conveniently pre-package smaller quantites of individually packaged items by, for example, size or color. In this manner less than case lot quantities of a particular size or color of an item may be economically packaged, or a case lot quantity containing a mixed assortment of items may be packed eliminating the necessity of case-lot purchases of each size or color.\nIn order to perform these functions, a distributor pack must provide an economical packaging medium for distribution of moderate quantities of individual pre-packaged items in order to minimize the additional packaging cost to the product manufacturer. In addition, such packaging must be sufficiently strong to provide adequate protection for the individual pre-packaged products contained therein, and to prevent loss or damage to the individual pre-packaged items or their package.\nThe present invention provides an economical distributor pack carton which does not require the addition of paperboard closure devices to close the bottom of the carton, or require the use of two-piece telescoping cartons as heretofore necessary. The distributor pack carton of the present invention has the bottom closure thereof formed from an integral portion of the carton body and protects individual packages within the carton against damage.\nThe carton body is formed as an open end tube which may be readily constructed using high speed gluing techniques since the tube is formed by single-line strip gluing along a single straight edge. In addition, the bottom closure of the open ended tube may be quickly and conveniently set up, to enable the distributor pack to be quickly filled with individual pre-packaged items at a minimum of time and expense. The open ended tube construction having the bottom closure formed from an integral portion of the tube body eliminates the necessity of complex packaging assembled about an array of individual consumer packaged goods, or expensive and complex operations requiring on-site gluing or the like.\nIn one embodiment of the invention, predetermined portions of the panels from which the open end tube carton is formed are removed to display a portion of the individual pre-packaged items to be contained within the pack. The provision of such openings in the distributor pack carton permits the ready determination of such things, for example, as size, quantity, color, price, date or \"UPC\" coding on the individual packages. In addition, providing these openings in the body of one carton allows two such distributor pack cartons to be nested together and formed from a single standard sheet of paperboard. Since standard sheets of paperboard are customarily formed in a square or rectangular configuration, individual carton blanks may be inter-nested or laid out on the standard sheet of paperboard in complementary form to provide maximum utilization of the paperboard in forming the carton blank. The formation of a top closure of one of the open ended tubes from the paperboard material removed to form the openings in the panels of the other complementary inter-nested carton blank permits the greatest efficiency and utilization of the paperboard stock from which these cartons are formed and substantially eliminates wasted paperboard."}
-{"text": "Compositions used to enhance cosmetic products are known in the art. Such compositions include those that are applied over top compositions such as lipstick to provide attributes such as gloss, lubricity and transfer-resistance of the cosmetic product they are applied over. These enhancement products utilize a variety of polymeric fluids and film forming technologies. For example, acrylic film-formers that are incorporated in lipstick overcoat products such as CSI Incorporated's \"Sealed with a Kiss\" are delivered in a volatile vehicle, alcohol, which is spread over the lipstick surface.\nAlternative topcoat products to those described above are disclosed in Japanese Pat. Application Number HEI 5 [1993]-221829, published Aug. 31, 1993. Said overcoats are reputed to exhibit improved durability of makeup effect, suppression of color transfer, and improved applicability. Said topcoats comprise from 0.2 to 25% of silica powder and/or alumina powder and from 75% to 99.8% of a perfluoropolyether of general formula: ##STR1## wherein R.sup.1 though R.sup.5 are independent fluorine atoms, perfluoroalkyl groups, or oxyperfluoroalkyl groups; the value of p, q, and r is at least zero; wherein the perfluoropolyether molecular weight is from about 500 to about 10,000, wherein P, Q and R may be equal, but, not zero. The preferred perfluoropolyether disclosed therein is a commercially available product known as Fomblin HC-04, HC-25, and HC-R available from Montefluosu of Milano, Italy.\nWhile such compositions may provide certain advantages, it has been found that they often disrupt the primary advantages of the cosmetic products they are applied over. For example, cosmetic products compromise their gloss or feel attributes in order to improve the long wear properties provided by the composition that is applied over top the cosmetic product. Alternately, cosmetic products must sacrifice long wear properties in order to improve the gloss and or feel attributes provided by the such compositions."}
-{"text": "As of now, in the quest to intensify the iron-making process the efforts are directed to improving the blast-furnace process through optimum distribution of a gas flow over a shaft section of the blast furnace, which determines the quality of iron produced, the specific fuel consumption for its production and the blast furnace output as a whole. The uniform distribution of the gas flow over the furnace section is attained by controlling the rate of blowing and changing the layer contours of ore and coke components of the burden which is fed into a throat zone by a charging device installed on the furnace cone. Therefore, the requirements that are placed upon the charging device are to provide a desired contour of the burden layer over the furnace cross section with minimum nonuniformity in the circumferential grain-size distribution of the burden in the throat zone, and, essentially, a controlled thickness of the burden layer charged.\nConventional blast furnace charging apparatuses include a double-cone charging apparatus equipped with movable throat plates, which is disclosed in DE A 2125062.\nThe above apparatus for charging a shaft furnace comprises a bin having an inlet and an outlet, the bin being connected through flanges with a top portion of the furnace. The bin outlet is located in a furnace chamber, within its throat zone, and closed by a cone-shaped locking member. The locking member is adapted to move vertically. The furnace throat zone is provided with movable plates which carry distributing plates with a reflecting surface facing the bin outlet. The plates are uniformly circumferentially disposed and adapted to radially move along the guides which are rigidly fixed to the furnace shell.\nIn operation of the apparatus, the locking member is displaced towards the furnace chamber, and a burden is admitted through an annular gap formed to the peripheral region of the throat zone. If the blast-furnace process is disturbed, it is required to change the contour of the charged burden layer. To this end, the plates are radially displaced towards the furnace axis until they reach a position in which the path of the falling burden intersects the surface of the distributing plates, which results in altering the burden flow direction so that the burden moves towards the furnace axis. The above prior art apparatus, however, changes the contour of the burden layer charged only in the peripheral annular region due to the restricted displacement of the movable plate, which reduces the potentials of the apparatus, limiting the burden laying control only to 1/3 the furnace radius.\nMost closely approaching the present invention by the combination of technical features and the achieved result is an apparatus for charging a shaft furnace, comprising a bin for a burden, the bin having an inlet and an outlet, a unit for distributing the burden over a cross section of the furnace, said unit being mounted in a throat zone beneath the bin outlet and adapted to rotate about the furnace axis. The burden distributing unit comprises a horizontal member having at least two guiding members circumferentially disposed around its periphery, each of the guiding members being connected with the horizontal member and consisting of two segments which are arranged sequentially in the direction of the burden flow. A first segment is located in the immediate vicinity to the horizontal member, and a second segment is comprised of at least two plane faces and has a burden shoot edge facing the furnace throat zone (WO Application No. 92/19776, C21B 7/20).\nOwing to the design of the distributing unit, the burden body leaving the bin outlet is divided into at least two flows which form a charge layer in the throat zone with approximately even circumferential grain-size distribution. This is attained due to laying the burden pellets of the same size symmetrically about the furnace axis, thereby significantly reducing the nonuniformity of the circumferential burden weight distribution in the furnace. The uniform circumferential arrangement of the guiding members, with the first segment being disposed in the immediate vicinity to the horizontal member, makes it possible to change the radial direction of the burden flow to the circumferential one. On the guiding member second segments, the circumferential direction of the burden flow changes to the tangential one or towards the furnace axis under the effect of the resultant of frictional, centrifugal and Coriolis forces. This enables the burden to be charged to any throat region and provides the possibility to form a required stockline contour in the throat zone. This is attained by a single type of motion, in particular, by rotation of the burden distributing unit about the furnace axis. The above structural features of the charging unit provide the improvement in the blast-furnace process parameters and the reduction in energy losses. The rotary drive of the unit has a simpler structure.\nIn operation, a burden portion is charged from the burden bin on the horizontal member of the burden distributing unit and further along the guiding members into the furnace chamber. A burden layer is formed in the furnace throat zone by virtue of rotation of the burden distributing unit. Owing to the possibility to vary the rotational speed of said unit, the apparatus provides a desired stockline contour in the shaft furnace. This improves the blast-furnace process characteristics. However, a required stockline contour is not provided when laying the burden by the above apparatus in the region adjacent to the furnace axis, since the non-optimal spatial orientation of the plane faces of the guiding member second segment hinders the admission of a sufficient amount of material into the furnace axial region. As a consequence, the optimum gas distribution over the furnace cross section is not provided, resulting in the increased coke consumption and some reduction in the output. In addition, in the prior art apparatus the position of plane faces of the guiding members fails to be unambiguously defined which makes the manufacture of the apparatus as such problematic, and further hampers the attainment of a desired stockline contour in the furnace throat zone."}
-{"text": "One type of offshore terminal for mooring a ship includes upper and lower beams extending from the bow of the ship and supporting a turret in rotation about a vertical axis relative to the ship. The lower portion of the turret is anchored to the sea floor, as by heavy chains. An upper bearing arrangement that mounts the turret to the upper beam, includes a thrust bearing that supports the weight of the turret and the load thereon, and that also forms a radial bearing that resists horizontal movement of the turret. A lower bearing arrangement that mounts the turret to the lower beam includes only a radial bearing. The turret is rigid, and the upper and lower bearings must be precisely aligned to avoid large bending stresses on the turret which would result in large radial loads on the bearings that limit their useful life.\nIn practice, it is found difficult to mount the radial bearings on the upper and lower beams to be precisely aligned. For example, even if the bearings are precisely aligned when the turret is installed on the vessel, misalignment will occur at the site where the terminal is installed when heavy chains are attached to the turret. The heavy weight of the chains causes the upper beam to bend, and in bending its outer ends pivot, which results in misalignment of the bearings. A comprehensive analysis can be made to predict the degree of bending of the upper beam, but in practice such analysis is very difficult to perform accurately. Also, it is difficult to install the turret at a shipyard with the required amount of misalignment and consequent bending load on the turret, which is predicted to be compensated for when the chains are installed. A bearing arrangement which avoided the need to accurately predict the amount of bearing misalignment during final installation and the need to install the bearings under high load with high precision, would be of considerable value."}
-{"text": "The present invention relates to an injection molding machine having an injection device and a mold opening and closing device, particularly to a motor-driven injection molding machine using motors as drive sources in the injection device and the mold opening and closing device.\nAs is well known, according to a motor-driven injection molding machine, respective functions of injection, mold opening and closing, measurement and ejector, are realized by individually installed motors. Induction motors are normally used for such motors. An explanation will be given of an example of a motor drive for such a motor in reference to FIG. 1. In FIG. 1, a motor drive 10 is for a three phase induction motor IM and includes an inverter section 20 and a rectifying circuit section 30. The inverter section 20 includes a plurality of switching elements 21 for carrying out switching operation for respective phases of three phases. The rectifying circuit section 30 is realized by a condenser input type diode rectifying circuit. The diode rectifying circuit includes a plurality of diodes 31 for carrying out rectifying operation for respective phases of a three phase power source PS.\nMeanwhile, the induction motor IM for the injection device or the mold opening and closing device generates regenerated electric power in deceleration thereof. Heretofore, in the motor drive, regenerated electric power of the induction motor IM is dissipated as heat by a resistor 41 provided to a dynamic brake circuit section 40. That is, when the induction motor IM generates regenerated electric power, the switching element 42 is switched on and the regenerated electric power is dissipated by the resistor 41. From a view point of reduction of electric energy, wasteful electric energy is dissipated. In addition thereto, there is constituted a heat radiating device for radiating generated heat of the resistor 41 and therefore, there poses a problem of an increase in dimensions.\nAs a measure with regard to the above-described regenerated electric power, there is provided even a system enabling power source regeneration by replacing the rectifying circuit section 30 by a power source regenerating converter. However, the power source regenerating converter is expensive.\nFurther, other than the above-described, it is the current situation that heretofore, according to the injection molding machine, no measure has been taken for electricity interruption except a necessary minimum power source backup function. Therefore, when molding operation is stopped by electricity interruption in the midst of a molding cycle during the operation, in the case in which resin is not sufficiently charged in a mold, an ejector mechanism cannot eject unmolded product in the mold. In this case, the unmolded product, that is, a failed molded product must be taken out by manual operation. Time and labor is required in taking it out and depending on cases, the mold must be disassembled. Further, when the operation is stopped in the midst of the molding cycle by electricity interruption, much time and labor is required and skill is needed in resetting peripheral apparatus of the injection molding machine in restarting thereof.\nHence, it is an object of the present invention to provide an injection molding machine capable of reducing power consumption by achieving effective utilization of regenerated electric power of a motor.\nIt is another object of the present invention to provide an injection molding machine capable of utilizing regenerated electric power of a motor to operation in electricity interruption.\nAn injection molding machine according to the present invention comprises a drive source comprising at least one motor and a motor drive, a charger for storing regenerated electric power of the motor and electric power from a power source, a charging and discharging circuit connected to the charger for charging electric power to the charger and for discharging electric power therefrom, and a controller for controlling the charger to charge electric power and for controlling the charger to supply electric power to the motor. The controller controls the charger by controlling the charging and discharging circuit based upon an operating condition of the injection molding machine.\nAccording to another aspect of the present invention, a method of providing drive power to an injection molding machine is provided. The method comprises the steps of controlling a charger to charge electric power from a power source and to supply electric power to a motor of the injection molding machine. A controller provides control signals to a charging and discharging circuit based upon an operating condition of the injection molding machine."}
-{"text": "In recent years, implementation of \u201ccloud-based\u201d services, high-performance computing (HPC) and other activities employing data centers and the like have seen widespread adoption. Under a typical data center installation, a large number of servers installed in server chasses and server racks are interconnected in communication using network links (e.g., Ethernet or InfiniBand) and various switching mechanisms, such as switch blades/modules and \u201ctop-of-rack\u201d (ToR) switches.\nUnder aspects of HPC, a very large number of compute nodes are implemented to solve various tasks in a parallel or substantially parallel manner. In essence, each compute node performs one part of a larger, more complex task. In order to implement this scheme, there needs to be input data to and output data must be exchanged among compute nodes. This data is communicated using one or more interconnects.\nVarious types of interconnects are used to interconnect the computer nodes in an interconnect hierarchy. For example, at the top of the hierarchy are interconnects between computing cores in the same processor. At the next level are interconnects between processors in the same platform, such as server blade or module. Next are interconnects between platforms, such as a backplane in a blade server. This is followed by interconnects between server chassis and/or ToR switches, interconnects between server racks, and finally interconnects between data centers. Generally, the communication bandwidth between nodes is reduced as one moves farther down the hierarchy.\nIn addition to latencies corresponding to transfers across the interconnect links themselves (which is a function of the link bandwidth and length and switching speed), significant latencies result from operations performed at the interfaces to the interconnects and/or additional processing/required to prepare the data for transfer across various interconnect links in the interconnect hierarchy. These data transfer latencies collectively reduce the communication performance and therefore the performance of the overall HPC implementation, and may represent a significant portion of the total latency (processing and data transfer) for the compute nodes.\nAnother important aspect of HPC is the software architecture. Implementing software using 10's of thousands of compute nodes in a parallel manner requires a significantly different software architecture than that employed for conventional applications. In addition, specific software modules have been developed for using corresponding types of interconnects, such as software modules used for communicating over InfiniBand."}
-{"text": "This invention relates to a process for the treatment of strongly acidic cation exchange resin catalysts based on styrene/divinyl benzene copolymerizates.\nRecently, cation exchange resin catalysts have been employed in processes for the ecologically beneficial implementation of acid-catalyzed syntheses. The syntheses in question are for instance esterifications, splitting of esters, hydrolyses, condensations, hydrations as well as alkylations and acetylations of aromatics. Unlike liquid acids, they have the advantage that the catalyst can be easily separated from the product and that no waste acid mixtures are produced as in the conventional homogeneous catalysis.\nA prerequisite for the practical use of a solid cation exchanger instead of a liquid acid is, in addition to sufficient selectivity and space/time yield, the thermal stability of the copolymerizates under the respective reaction conditions.\nStrongly acidic styrene/divinyl benzene copolymerizates, the core of which has been substituted with halogen, are particularly thermally stable and are used at temperatures in the range of 100.degree. to 200.degree. C. for acid-catalyzed syntheses, such as for the hydration of lower olefins or for alkylation reactions."}
-{"text": "Automotive trim parts, particularly body side moldings of various sorts have been known for many years. Such trim parts have both decorative and protective functions. Typically, there are at least three pieces horizontally aligned and attached on the side of a vehicle: a piece for the front fender; a piece for each door; and a piece for the rear fender.\nOne known type of molding is made of a polymer coextruded onto a metal carrier strip. The metal is typically aluminum, stainless steel or bimetal and a portion of the metal may be left exposed to enhance the appearance of the molding. The underlying metal carrier generally has an outwardly convex cross-section and its longitudinal edges are bent inwardly to form lips running the length of the piece. This arrangement thus yields a hollow molding which is generally fastened to a vehicle by clipping of the lips to the vehicle body. This type of hollow molding requires a member to be fastened at each of its ends. Such end members close off the hollow ends, precluding each end from undesirably snagging other objects and further enhancing the appearance of the molding by concealing section components. The end members are generally of polymer material and are often differently shaped from each other. For example, the design of certain automobiles requires the front end of a front door molding to be tapered so that it does not catch on the rear edge the of front fender when the door is opened. The rear end of a rear fender molding is often shaped to match the shape of the rear wheel well. In any case, this arrangement results in a trim piece made up of the extruded strip and end members, which must be assembled and fastened to each other during manufacture, and having aesthetically undesirable parting lines between the strip and end members. Its hollow construction is considered advantageous however, for its economical use of materials and light-weight aspect.\nAnother known type of molding comprises a extruded piece of solid plastic, possibly having a decorative metal strip coextruded as part of its exterior face. The molding is usually fastened to an automobile using an adhesive such as double-sided sticky tape. The extruded strip may be cut to obtain a blunt end during manufacture, but if section components are to be concealed, end members are required as with the hollow-type molding described above. Further, if a piece with shaped ends is desired, end members must generally be used because cutting of the extruded plastic to obtain, for example, a tapered shape results in a cut plastic surface exposed as part of the decorative face of the molding, and this is generally unacceptable in the automotive industry.\nA trim piece of hollow construction, not requiring end members to be attached to its ends and which may be readily attached to a vehicle without the use of clips would be considered an improvement over these known types of moldings."}
-{"text": "Portable electronic devices, such as laptops, tablets, and smartphones, are increasingly used in settings in which sensitive information may be stored on the device. For example, a user may send and receive corporate email on their smartphone or may access, edit, and save internal corporate documents on a tablet device, such as via a remote network connection. In addition, it is becoming increasingly common for users to supply their own devices for use at work under policies commonly referred to as \u201cbring your device\u201d or BYOD policies. However, devices used under such BYOD policies present security risks for information technology (\u201cIT\u201d) departments, who traditionally provided corporate devices to users with particular security features enabled."}
-{"text": "1. Field of the Disclosure\nThe present disclosure relates generally to a multiple-input multiple-output (MIMO) technology for an antenna structure, and more particularly to an electronic device including the antenna structure.\n2. Description of the Related Art\nThe advance of mobile communication and antenna technologies are accompanied by an increase of the number of antennas included in electronic devices. However, as electronic devices are generally small-sized, the available spaces for mounting the antennas is limited.\nFurther, because antennas are normally mounted in the limited interior spaces of the electronic devices, the antennas are often arranged very close to each other. Consequently, these neighboring antennas may interfere with each other, for example, when signals of adjacent frequencies are transmitted and received, thereby lowering of the performance of the antennas.\nAdditionally, when a housing of an electronic device is formed of a conductive material (for example, metal), the radiation performance of the antenna may also be affected."}
-{"text": "Recently, mobile communication devices and small personal computers are mounted with a camera module capable of generating a digital image or a digital video corresponding to an outside light. A camera module includes an image sensor module changing an outside light to an image and a lens focusing the outside light to the image sensor module. A conventional camera module is disadvantageously fixed by the lens and the image sensor module, and a distance between the lens and the image sensor module cannot be adjusted to make it difficult to obtain a high quality image desired by a user.\nRecently, a distance between the lens of a camera module and the image sensor can be adjusted by a VCM (Voice Coil Motor) to obtain a very high quality image from the camera module. A conventional VCM is configured in such a manner that an image sensor module is arranged at a rear surface, or a rear side of a base, and an IR (Infrared) filter is arranged at the rear surface of the base.\nHowever, the IR filter recently suffers from a disadvantage in that, concomitant with increased resolution of an image sensor module, quality of an image or a video generated by the image sensor module greatly deteriorates in a case a foreign object such as dust is penetrated into the IR filter."}
-{"text": "1. Field\nExample embodiments disclosed herein relate to a display device to provide a three dimensional (3D) image and an internal image processing method, and more particularly, to an image processing apparatus and method that may decrease a visual fatigue when viewing a 3D image.\n2. Description of the Related Art\nRecently, interest in three dimensional (3D) image processing has increased. A 3D image may be configured by providing images corresponding to different views with respect to a plurality of views. For example, the 3D image may be a multi-view image corresponding to a plurality of views, or a stereoscopic image that may provide a left eye image and a right eye image corresponding to two views.\nCompared to a two dimensional (2D) display which displays only an existing plane image, a 3D display may provide a relatively higher sense of reality to a user. However, the 3D display may express different image information with respect to a left field of view and a right field of view and thus, may cause increased visual fatigue as compared to the 2D display. Accordingly, there is a desire for a solution to decrease visual fatigue.\nConventional methods for decreasing visual fatigue include a method of uniformly decreasing depth information, a method of increasing a contrast in the case of a near side region and decreasing the contrast in the case of a far side region, and the like."}
-{"text": "Social networking has become a popular medium for online collaboration. Maturing Web 2.0 technologies have influenced competing social networking sites (for example, Orkut, Facebook, etc.) to look at ways to provide innovative features. Also, telecommunications (telecom) technologies have experienced changes such as, for example, content and capabilities owned by the operator are now available for developers to combine. Additionally, social networking now includes the use of mobile devices to enable communications with on-line social networking buddies.\nHowever, traditionally, social networks have worked on relatively static and/or slow-changing information, that is, defined once by a user and changed infrequently. Groups that are defined in a social network in existing approaches are similarly static and pre-defined in terms of membership. Further, ways of communicating with group members are also statically defined (for example, leave a scrap, write on a wall, etc.)."}
-{"text": "Schizophrenia is a chronic debilitating mental illness affecting about one percent of the population. The disease manifests in delusional behavior, dysfunctional thinking, agitated body movement, social withdrawal, and depression. Schizophrenia patients suffer a profoundly reduced quality of life, and are ten times more likely to commit suicide that the general population.\nDopamine (particularly D2 and D3) antagonists are well recognized as improving symptoms of schizophrenia, and have been used clinically as such for decades. In the past twenty years it has become recognized that treatment of schizophrenia, as with many mental illnesses, benefits from engaging multiple receptors including serotonergic and adrenergic. Despite, literally, dozens of approved drugs to treat schizophrenia the disease remains poorly treated in many patients. Side effects of current medications include: dyskinesia, akathisia, weight gain, mood disturbances, sexual dysfunction, sedation, orthostatic hypotension, hypersalivation, and (in some cases) arganulocytosis.\nAmisulpride (4-amino-N-(((1-ethyl-2-pyrrolidinyl)methyl)-5-(ethylsulfonyl))-2-methoxybenzamide) is an antipsychotic patented in 1981. Amisulpride binds selectively to the human dopaminergic D2 (Ki 2.8 nM) and D3 (Ki 3.2 nM) receptor subtypes without any affinity for D1, D4 and D5 receptor subtypes. Unlike classical and atypical neuroleptics, amisulpride displays low affinity for serotonin, alpha-adrenergic, histamine receptor subtypes, muscarinic receptors and sigma sites though it has also been demonstrated to bind 5-HT2B and HT7a receptors with low double digit nM Ki. This ability of amisulpride to bind 5-HT receptors is thought to result in amisulpride's ability to treat symptoms of depression (sometimes noted in schizophrenia patients). Interestingly, compared to other antipsychotics, amisulpride is not noted to have any activity at the 5-HT2a receptor.\nDespite the unique activities of amisulpride, amisulpride has low ability to cross blood brain barrier (BBB) to interact with the receptors in the brain. In a 2014 study, passive diffusion of amisulpride across a PAMPA membrane (Pe) was the lowest of 30 psychiatric drugs tested. Thus, dosing of amisulpride is high, typically 400 to 800 mg/d (though up to 1,200 mg/day is not uncommon). Such a high dose may cause adverse effects to the treated subjects."}
-{"text": "1. Field of the Invention\nThis invention relates to a light shielding mechanism for a quadrilateral opening and more particularly to a light shield mechanism for such an opening within a single lens reflex camera and the like.\n2. Description of the Prior Art\nIn single lens reflex cameras, provided with a mechanism for automatically determining the aperture size or the extent of the exposure time or similar parameters in accordance with photographing conditions such as the brightness of the field to be photographed, the camera employs an internal light receiving photometric system in which a photosensitive portion for measuring the light within the field to be photographed arranged within the region of the light path emanating from the field to be photographed and transmited through an image lens in order to determine the aforementioned parameters for film exposure. However, it is not possible to disregard the influence of counter-incident light which falls on the light path from the eyepiece portion of the view finder optical system. Particularly where the eye piece portion is exposed when the photograph is taken by means of a self-timer or remote control mechanism, a relatively large error is produced by the presence of the counter-incident light acting on the light being measured relative to the field being photographed.\nFor this reason, it is often required that a mechanism be provided for shielding the view finder eyepiece in single lens reflex cameras and the like, and it is further necessary to add such a mechanism for shielding a light projection optical system when a slide change is being effected within a slide projector to thereby enhance projection effects, and similarly it is often required that such a shield mechanism be provided within optical illumination systems in telescopic photograhic apparatus or the like.\nIt is, therefore, an object of the present invention to provide a light shield mechanism for a camera view finder eyepiece and the like which reduces the area of the shield blades incorporated within said mechanism where the optical path across which the light shield is employed comprises a quadrilateral opening such as a rectangle or the like.\nIt is a further object of the present invention to provide an improved shield mechanism which minimizes the space required for the shield blades within the periphery about the quadrilateral light path opening and to provide a light sield mechanism which is both simple in construction and which readily attains the desired light shielding function."}
-{"text": "1. Field of the Invention\nThe present disclosure relates to methods of treating light-colored wood substrates to produce a wood substrate with a low gloss, open pore appearance that provides enhanced protection from ultraviolet (\u201cUV\u201d) radiation and a \u201creal\u201d wood look. The present disclosure also relates to the wood components produced by the methods disclosed herein.\n2. Description of the Related Art\nReal wood trim pieces with very glossy finishes in some automotive vehicles are difficult to distinguish from plastic woodgrain films. A simple polyurethane (herein \u201cPU\u201d) coating can have the desired appearance but typically cannot meet rigorous performance targets for, for example, long term resistance to fading upon extended exposure to sunlight.\nHowever, real wood remains a popular choice for the interior of upscale and luxury automobiles. In order to meet rigorous performance targets multiple coats of high gloss finishes are routinely used with real wood trim pieces in automotive vehicles in an attempt to meet the light resistance standard and other performance criteria, but the resulting high gloss finish makes the real wood appear to the ordinary consumer to be less desirable woodgrain plastic films.\nA method of achieving the necessary light resistance properties while maintaining a \u201creal\u201d wood appearance, that is, not a high gloss finish, is of interest."}
-{"text": "A lens formed of a high-refractivity low-dispersion glass enables the downsizing of an optical system while correcting a chromatic aberration when combined with a lens formed of a high-refractivity high-dispersion glass. It hence occupies an important place as an optical element for constituting an image-sensing system or a projection optical system such as a projector.\nJP 2007-269584A discloses such a high-refractivity low-dispersion glass. The glass disclosed in JP 2007-269584A has a refractive index nd of 1.75 to 2.00 and has a Ta2O5 content in the range of 0 to 25 mass %, and all of the glasses that have a refractive index nd of at least 1.85 contain a large amount of Ta2O5. That is because the introduction of a large amount of Ta2O5 is indispensable for securing glass stability in the region of high refractivity such as a refractive index nd of 1.75 or more. For such a high-refractivity low-dispersion glass, Ta2O5 is a main and essential component.\nMeanwhile, tantalum (Ta) is an element having a high rarity value and is in itself a very expensive substance. Moreover, rare metal prices are recently soaring worldwide, and the supply of tantalum is deficient. In the field of glass production, tantalum as a raw material is deficient, and if such a situation continues, it may be no longer possible to maintain the stable supply of high-refractivity low-dispersion glasses that are essential and indispensable in the industry of optical apparatuses."}
-{"text": "1. Field of the Invention\nThis invention relates generally to the field of semiconductor device manufacturing and, more particularly, to a method and apparatus for adaptively scheduling tool maintenance.\n2. Description of the Related Art\nThere is a constant drive within the semiconductor industry to increase the quality, reliability and throughput of integrated circuit devices, e.g., microprocessors, memory devices, and the like. This drive is fueled by consumer demands for higher quality computers and electronic devices that operate more reliably. These demands have resulted in a continual improvement in the manufacture of semiconductor devices, e.g., transistors, as well as in the manufacture of integrated circuit devices incorporating such transistors. Additionally, reducing the defects in the manufacture of the components of a typical transistor also lowers the overall cost per transistor as well as the cost of integrated circuit devices incorporating such transistors.\nGenerally, a set of processing steps is performed on a lot of wafers using a variety of processing tools, including photolithography steppers, etch tools, deposition tools, polishing tools, rapid thermal processing tools, implantation tools, etc. The technologies underlying semiconductor processing tools have attracted increased attention over the last several years, resulting in substantial refinements. However, despite the advances made in this area, many of the processing tools that are currently commercially available suffer certain deficiencies. In particular, such tools often lack advanced process data monitoring capabilities, such as the ability to provide historical parametric data in a user-friendly format, as well as event logging, real-time graphical display of both current processing parameters and the processing parameters of the entire run, and remote, i.e., local site and worldwide, monitoring. These deficiencies can engender nonoptimal control of critical processing parameters, such as throughput, accuracy, stability and repeatability, processing temperatures, mechanical tool parameters, and the like. This variability manifests itself as within-run disparities, run-to-run disparities and tool-to-tool disparities that can propagate into deviations in product quality and performance, whereas an ideal monitoring and diagnostics system for such tools would provide a means of monitoring this variability, as well as providing means for optimizing control of critical parameters.\nOne technique for improving the operation of semiconductor processing line includes using a factory wide control system to automatically control the operation of the various processing tools. The manufacturing tools communicate with a manufacturing framework or a network of processing modules. Each manufacturing tool is generally connected to an equipment interface. The equipment interface is connected to a machine interface which facilitates communications between the manufacturing tool and the manufacturing framework. The machine interface can generally be part of an advanced process control (APC) system. The APC system initiates a control script based upon a manufacturing model, which can be a software program that automatically retrieves the data needed to execute a manufacturing process. Often, semiconductor devices are staged through multiple manufacturing tools for multiple processes, generating data relating to the quality of the processed semiconductor devices.\nVarious tools in the processing line are controlled in accordance with performance models to reduce processing variation. Commonly controlled tools include photolithography steppers, polishing tools, etching tools, and deposition tools. Pre-processing and/or post-processing metrology data is supplied to process controllers for the tools. Operating recipe parameters, such as processing time, are calculated by the process controllers based on the performance model and the metrology information to attempt to achieve post-processing results as close to a target value as possible. Reducing variation in this manner leads to increased throughput, reduced cost, higher device performance, etc., all of which equate to increased profitability.\nCommonly, a processing tool undergoes periodic preventative maintenance procedures or calibrations to keep the tool in optimum operating condition. For example, polishing tools include polishing pads that are periodically conditioned or replaced. Etch tools and deposition tools are periodically cleaned using both in situ cleans or complete disassembly cleans. Steppers are periodically calibrated to maintain alignment accuracy and exposure dose consistency. Many of these preventative maintenance procedures are performed at discrete intervals based on vendor recommendations, past history, and expected degradation rates of consumable items used in the tools. The use of fixed preventative maintenance intervals is not always an effective solution for optimizing tool and line efficiency. If the maintenance activities are performed more often than actually needed, the efficiency of the line and the operation cost of the tool is increased. If maintenance activities are performed less often than needed, product quality and tool reliability may be degraded.\nAnother potential drawback of scheduled maintenance intervals is that unexpected conditions may arise during the time period between maintenance activities. For example, the calibration of a tool (e.g., focus on a stepper) may drift outside expected tolerances. A process controller used to control a particular tool may adjust the operating recipe of the tool trying to correct for what it thinks is normal process variation, while in actuality, the variation is caused by the unexpected condition. Over time, the process controller may adjust its performance model based on post-processing metrology feedback. In some situations, the process controller might not be able to stabilize the process. These control problems may result in increased variation or even defective wafers.\nThe present invention is directed to overcoming, or at least reducing the effects of, one or more of the problems set forth above.\nOne aspect of the present invention is seen in a method for adaptively scheduling tool maintenance. The method includes controlling an operating recipe of a tool using a plurality of control actions, monitoring the control actions to identify a degraded tool condition, and automatically initiating a tool maintenance recommendation in response to identifying the degraded tool condition.\nAnother aspect of the present invention is seen in a manufacturing system including a tool, a process controller, and a tool health monitor. The tool is adapted to process a workpiece in accordance with an operating recipe. The process controller is adapted to control the operating recipe of the tool using a plurality of control actions. The tool health monitor is adapted to monitor the control actions to identify a degraded tool condition and automatically initiate a tool maintenance recommendation in response to identifying the degraded tool condition."}
-{"text": "This invention relates to a connector plug for terminating polarization maintaining (PM) optical fibers.\nIn optical fiber communications, connectors for joining fiber segments at their ends, or for connecting optical fiber cables to active or passive devices, are an essential component of virtually any optical fiber system. The connector or connectors, in joining fiber ends, for example, has, as its primary function, the maintenance of the ends in a butting relationship such that the core of one of the fibers is axially aligned with the core of the other fiber so as to maximize light transmissions from one fiber to the other, or, put another way, to reduce insertion loss. Another goal is to minimize back reflections. Alignment of these small diameter fibers is extremely difficult to achieve, which is understandable when it is recognized that the mode field diameter MFR of, for example, a singlemode fiber is approximately nine (9) microns (0.009 mm). The MFR is slightly larger than the core diameter. Good alignment (low insertion loss) of the fiber ends is a function of the transverse offset, angular alignment, the width of the gap (if any) between the fiber ends, and the surface condition of the fiber ends, all of which, in turn, are inherent in the particular connector design. The connector must also provide stability and junction protection and thus it must minimize thermal and mechanical movement effects.\nIn the present day state of the art, there are numerous, different, connector designs in use for achieving low insertion loss and stability. In most of these designs, a pair of ferrules (one in each connector), each containing an optical fiber end, are butted together end to end and light travels across the junction. Zero insertion loss requires that the fibers in the ferrules be exactly aligned, a condition that, given the necessity of manufacturing tolerances and cost considerations, is virtually impossible to achieve, except by fortuitous accident. As a consequence, most connectors are designed to achieve a useful, preferably predictable, degree of alignment, some misalignment being acceptable.\nHowever, in connecting or terminating polarization maintaining (PM) fibers, such is not the case. Many optical fiber components, such as, for example, interferometers and sensors, lasers, and electro-optic modulators, are extremely sensitive to and dependent upon, for proper operation, the polarization of the light. Even very slight alterations or changes in the light polarization orientation can result in wide swings in the accuracy of response of such devices. PM fiber has polarization-dependent refractive indices, and the speed of light in an optical fiber is inversely proportional to the magnitude of the refractive index. A birefringent optical fiber is one having two polarizations having different velocities of propagation, thus giving rise to a xe2x80x9cfastxe2x80x9d wave and a xe2x80x9cslowxe2x80x9d wave. In a PM fiber, the polarization of a linearly polarized light wave input to the fiber, with the direction of polarization parallel to that of the one of the two principal polarizations, will remain or be maintained in that polarization as it propagates along the fiber, hence the term xe2x80x9cpolarization maintaining.xe2x80x9d If the polarization of the light wave is to be maintained at a splice or other connection, the principal axes of birefringence of the two joined fibers must be aligned in parallel, otherwise there will be polarization cross-coupling, i.e., crosstalk, which is highly undesirable. Thus, where two PM fibers, for example, are to be connected together, they should be terminated carefully to reduce the crosstalk during the connectorization process. Also, the connectors must be capable of aligning then maintaining the fiber orientation to the connector key position. Connectors with tolerances adequate for connecting non-PM fibers usually are inadequate for maintaining polarization alignment at the connector junction.\nTypical PM connector requirements are an insertion loss of less than 0.3 dB, and the prior art PM connector arrangements comprise numerous, different connector configurations aimed at meeting these requirements for different connectors, such as an SC type connector as shown in U.S. Pat. No. 5,216,733 of Ryo Nagase et al. The connector of that patent comprises a ferrule body and a ring shaped flange having a keyway mounted on the periphery of the ferrule body. Alignment is achieved by rotating the ferrule body with respect to the flange keyway. The combination of ferrule and flange comprises a plug which is inserted into a push-pull SC connector having a key therein for mating with the flange keyway and springs bias the flange in the longitudinal direction to maintain the alignment.\nIn U.S. Pat. No. 4,784,458 of Horowitz, a splice joint for PM fibers is shown wherein aligned fibers are joined with UV curing epoxy, and the joint is overlaid with epoxy cement for rigidity. Such ajoint is permanent, and does not function as a connect-disconnect optical fiber connector.\nU.S. Pat. No. 5,561,726 of Yao discloses an apparatus for controlling the polarization state of the light within a fiber by squeezing a portion of the fiber to produce a birefringent fiber, and the squeezer is then rotated to change the polarization of the light within the fiber. The device is not a connector, but is intended for use with polarization sensitive devices such as interferometers and electro-optic modulators, however, it may also be used with connectors for connecting two PM fibers.\nIt is common practice in the prior art for creating PM fibers to include a pair of rods in the fiber cladding which extend parallel to the core as shown in U.S. Pat. No. 4,515,436 of Howard et al. Such rods, which are preferably of glass, are, in manufacture of the fiber, included in the fiber preform from which the fiber is drawn. As the fiber is drawn, the rods are accordingly diminished in diameter and are located within the cladding, preferably on either side of the core. The rods have different thermal expansion characteristics than the surrounding glass, and the stress they exert on the core causes the index of refraction to change along that axis. The axes then have different indices of refraction value and thus propagate light at different speeds. Variations on the two rod arrangement are also known, such as the elliptical stress member disclosed in U.S. Pat. No. 5,488,683 of Michal et al. Also, squeezing the fiber to create birefringence, as shown in the aforementioned Yao patent is feasible. The two rod PM fiber, so called xe2x80x9cPandaxe2x80x9d type PM fiber, however, has proven quite satisfactory in use, and it is toward the connectorization of such a fiber that the present invention is directed, although other types of PM fibers may be used with the present invention.\nIn the copending U.S. patent application Ser. Nos. 10/151,450 and 10/151,130 are shown, respectively, an apparatus and methods of tuning a PM connector plug and an adapter for the connector plug of the present invention the principles of which are applicable to any of a large number of optical fiber connectors, but are embodied in a modified LC connector in those applications. For optimum performance, i.e., maximum transmission of a polarized beam, it is highly desirable to provide accurate rotational positioning of better than xc2x11xc2x0 or even as accurate as  less than xc2xcxc2x0 between connectors equipped with polarization maintaining fibers.\nThe present invention is a connector plug for PM connectors which is intended for use with the adapter and tuning method of those Lampert et al. applications to achieve this desideratum. When a PM jumper cable, for example, is terminated by connectors, it is most desirable that the cable/connector combination be tuned to align the fiber slow axis with a reference point such as the connector key which may be the connector latching arm. In accordance with the present invention, there is provided a connector plug which is tunable to yield extremely accurate rotational positioning of the connector ferrules, or, more specifically, the fibers contained therein.\nThe connector plug of the invention comprises in a first embodiment, a tunable barrel assembly as shown in U.S. Pat. No. 5,481,634 of Anderson et al., the disclosure of which is incorporated by reference herein, a connector plug that includes the barrel-ferrule assembly for holding the end of an optical fiber extending axially therethrough and a housing for the assembly. A coil spring member contained within the housing surrounds the barrel, which is of tubular configuration, and bears against an interior wall of the housing and an enlarged flange member on the barrel, thereby supplying forward bias to the barrel-ferrule assembly relative to the housing. As is shown in the aforementioned U.S. patent applications, the barrel-ferrule assembly has an enlarged flange member which is hexagonal 1 in shape and has a tapered or chamfered leading surface that may be slotted. The housing, in turn, has a hexagonally shaped cavity, which provides any of six rotational positions for the flange and a tapered seating surface for the tapered surface of the flange. The dimensions of the cavity are such that the hexagonal barrel flange floats within the hexagonal cavity, in the Anderson et al. arrangement and can rotate about xc2x112xc2x0, which diminishes the tuning accuracy. In the PM connector of the Lampert et al. application Ser. No. 09/811074, this float is greatly reduced by dimensioning the hexagonal barrel flange, or nut, so that it is a virtual slip fit within the hexagonal cavity. Additionally, the flange is affixed to the barrel, hence, the barrel has only six positions, which are subject to the diminished float.\nIn the aforementioned U.S. patent application Ser. No. 09/811074, of Lampert et al, there is also shown a tunable nut arrangement wherein a hexagonal nut is a xe2x80x9clightxe2x80x9d press fit on the ferrule containing tubular portion of the barrel assembly, so that the nut may be rotated relative to the ferrule in extremely small increments, such as fractions of a degree. In the arrangement of that application, the nut may remain seated in the hexagonal seat in the connector plug body, and the ferrule rotated with respect thereto. As explained in that application, the xe2x80x9clightxe2x80x9d press fit is sufficient to prevent accidental or unwanted relative rotation after tuning. After tuning the nut (or flange) may be cemented in place. Such an arrangement is thus capable of producing the small increments of relative rotation that may be required for optimum PM tuning, and is a preferred arrangement for the first embodiment of the invention.\nIn accordance with the present invention, the latch of the plug, which includes a latching arm having latching shoulders on either side, preferably is of plastic material and has tapered sides, thereby imparting to the arm a cross-section shaped as a truncated wedge. The latching shoulders may also have tapered sides so that both the arm and the shoulders, alone or together have a taper. The adapter for use with the invention has a longitudinal slot therein which has sloping sides for receiving the tapered latching arm as shown in application Ser. No. 10/151,130. As is the practice with the latching arm, it is sufficiently resilient or elastic to snap into the slot once the latching lugs are engaged. In prior art arms and slot arrangements the slots and the arm both have straight sides, hence the arm is made slightly narrower than the slot to insure engagement, which gives rise to some play between the plug and the adapter which can alter the position of the polarization vector. The tapered arm and slot of the invention create a tight fit and eliminate such play so that the plug is firmly held in its latched position. In addition, the tapered engagement centers the plug relative to the adapter so that the ferrule/barrel assembly is properly aligned within the adapter.\nFurther in accordance with a second embodiment of the invention, the plug body has a rotatable extender cap therein which has longitudinally extending resilient arms, the distal ends of which bear against the hexagonal portion of the barrel assembly as a three jawed collet further to ensure alignment of the barrel and to eliminate virtually any float thereof. The barrel of the connector plug is an assembly comprising a tubular member which may have a notched end for tuning and an enlarged hexagonal nut affixed thereto on the outer surface of the tubular member and which has a tapered front section which finctions as a guide in conjunction with a tapered recess in the housing. As described hereinbefore, such an assembly is shown and described in the U.S. patent application Ser. No. 09/811074, the disclosure of which is incorporated herein by reference. Tuning of the connector is achieved by rotation of the extender cap which firmly holds the nut on the barrel. Hexagonal seat in the plug body is illuminated, and hence, the ferrule contained in the barrel can be incrementally rotated to produce, where necessary, incremental changes in the angular orientation.\nIn both embodiments of the invention, the latching arm has a tapered cross section as described hereinbefore, and the fitting of the parts to each other approaches a slip fit, so that float is minimized."}
-{"text": "1. Field of the Invention\nThe present invention relates to tubular connections for application to heavy wall, high performance casing in high pressure, critical service oil and gas wells. More particularly, the present invention relates to such a connection having an improved taper and buttress type threads.\n2. Description of the Prior Art\nIn recent years, oil and gas wells that have been drilled for exploration and production have commonly extended to depths in the range of 15,000 to 35,000 feet, where pressures and temperatures may exceed 15,000 psi and 250.degree. F. respectively. These conditions not only require tubulars (commonly known as Oil Country Tubular Goods (OCTG)) having higher strength, higher reliability and better resistance to failure under higher loads and corrosive applications, such as are disclosed in commonly assigned U.S. Pat. No. 4,354,882, but also require tubular connections which have higher strength and better resistance to severe stress applications. There are three basic types of OCTG's, each of which makes its own contribution to the drilling and completion of oil and gas wells. These types of tubulars are drill pipe, casing and tubing. This invention is primarily intended as a superior connection for heavy wall casing.\nTypically, the lengths of tubulars are threaded at each end and are joined together by an internally threaded coupling. The threads of both the tubes and the couplings must be able to carry the weight of the tubulars and couplings extending below it into the well and must also be able to withstand the high internal and external pressure encountered during drilling and producing of the wells. Since the casing commonly in use will have a weight of from 10 to in excess of 100 pounds per foot, the tensile loads which the threads in the uppermost casing connection must be able to withstand when there are 15,000 to 35,000 feet of tubulars and connections below it are tremendous.\nWith sufficiently high tensile loads, connections may fail due to rupture of the base metal within the connection or by sudden disengagement of the mated male and female threads. The latter of these failures is termed the \"pull-out\" mode of failure.\nThe casing connection must, of course, be leak-proof and a means of sealing must be provided. The threaded connections must be capable of being made up and disassembled without galling and without the danger of stripping threads or cross-threading. It is desirable after the joint is \"stabbed\" that it can be made up with a relatively small number of turns.\nTwo types of casing connections are in common use, coupled and integral joint. Coupled connections employ male threads on each end of the tube which are joined by a female threaded connector called a coupling. Tubes with male threads on one end and female threads on the other are referred to as integral joints since no couplings are required.\nThe following patents and publications illustrate the current state of the art.\nEaton U.S. Pat. No. 2,062,407 discloses a threaded pipe joint of the integral type which employs a buttress type thread having broad crests and roots and which is disclosed to have a reduced galling tendency (page 6, column 2, lines 48-50); easier and faster make-up (page 7, column 2, lines 28-40); and improved hoop strength at the end of the female section of the joint (page 3, column 2, lines 36-46).\nFrame U.S. Pat. No. 2,204,754 discloses a threaded pipe joint of the integral type having a modified acme type thread wherein the thread crests and roots are flat and parallel to the axis of the joint (page 2, column 1, lines 34-42). Five threads per inch are provided and are tapered at both ends, but not in the middle of the thread length, between 3/4\" and 13/4\" per foot (page 2, column 2, lines 21-27).\nRollins U.S. Pat. No. 2,885,225 discloses a drill pipe coupling employing a double thread having three to five threads per inch (column 2, lines 37-43) which is tapered about 11/2\" per foot (column 1, lines 63-67). This taper is said to result in easier and rapid make-up and to promote optimum stabbing (column 1, lines 65-67). The thread is provided with flat crests and roots which are parallel to the taper (column 4, lines 11-20).\nKloesel, Jr., et al., U.S. Pat. No. 3,355,192 discloses a pipe joint of the integral type having a modified buttress type double thread which is tapered 11/2\" per foot (column 3, lines 40-47, and see column 5, lines 40-45).\nYocum U.S. Pat. No. 3,346,278 discloses a tapered V-thread for pipe joints of the integral type wherein the external threads have a basic taper of between 20 percent and 80 percent of the basic taper of the internal threads (column 2, lines 25-28). This arrangement is said to produce a seal between the leading external threads and the innermost internal threads irrespective of whether either thread is at the extreme limits of the commercial taper tolerance (column 3, lines 34-40).\nNowosadko U.S. Pat. No. 3,427,707 discloses a non-threaded soldered pipe joint wherein the female member is machined to a larger taper than the male member so as to assure a seal at the leading end of the male member (column 2, lines 37-41).\nBlose et al. U.S. Pat. No. 3,224,799 discloses a threaded tube joint of the integral type which employs a buttress type thread tapered about 1\" per foot to permit deeper stabbing and more rapid make-up (column 2, lines 51-54; see also column 3, lines 44-47). U.S. Pat. No. 3,224,799 also suggests that the tolerance be such that the actual taper of the male threads will tend to be slightly deeper than that of the female threads (column 3, lines 47-55).\nBlount et al. U.S. Pat. No. 3,307,860 discloses a pipe coupling having an internally thickened central region and shoulders against which a resilient sealing ring is compressed by a shoulder on the ends of the pipe (column 6, lines 36-51; see also column 4, lines 49-55). The threads are of a tapered modified buttress type (column 3, lines 13-15).\nBlose et al. U.S. Pat. No. 3,572,777 discloses a pipe coupling thickened in the central region and provided with an internal groove for a resilient sealing ring which is compressed by the ends of the pipe (column 3, lines 1-18). The thread is of the acme or buttress type and is provided with a relatively steep taper (column 1, lines 5-16).\nAmerican Petroleum Institute (API) Standard 5B (March, 1979) discloses a tubular connection having a tapered buttress casing thread. The API tubular connection consists of a male joining element on the end of the pipe and a coupling member provided with cooperating female joining elements. The threads are on a taper of 3/4\" per foot on the diameter on sizes 133/8\" OD and smaller. The male joining element taper has a tolerance of from +0.0035 inch per inch to -0.0015 inch per inch. The female joining element taper has a tolerance of from +0.0045 inch per inch to -0.0025 inch per inch. The maximum disadvantageous taper spread (i.e., minus on coupling, plus on the pipe end) is 0.006 inch per inch. The API buttress threads have a substantially rectangular form with rounded corners and parallel crests and roots, a 0.200 inch thread pitch, a 3.degree. load flank, 10.degree. trailing flank, and a 0.062 inch thread height with a tolerance of plus or minus 0.001 inch.\nThe thread crests on the API male joining element threads have a length of 0.100 inch with a tolerance of from +0.000 inch to -0.003 inch (+0.000 inch to -0.005 inch on sizes 85/8\" and larger). The rounded corner on the load flank end of the crest has a radius of 0.008 inch with a tolerance of +0.002 to -0.000 and the rounded corner on the trailing flank end of the crest has a radius of 0.030 with a tolerance of from +0.002 to -0.000 inch. The roots have a length of 0.100 inch with a tolerance of from +0.003 to -0.000 inch (+0.005 to -0.000 on sizes 85/8\" or larger). The rounded corners on either end of the root have a radius of 0.008 inch with a tolerance of from +0.000 inch to -0.002 inch.\nThe crests of the API female joining element threads have a length of 0.099 inch with a tolerance of from 0.000 inch to -0.003 inch. The rounded corners at each end of the crest have a radius of 0.008 inch with a tolerance of from +0.002 inch to -0.000 inch. The roots have a length of 0.101 inch with a tolerance of +0.003 inch to -0.000 inch. The rounded corner of the load flank end of the root has a radius of from 0.008 inch with a tolerance of +0.000 inch to -0.002 inch. The rounded corner of the trailing flank end of the root has a radius of 0.030 inch with a tolerance of from +0.000 inch to -0.002 inch.\nThese API buttress thread tolerances result in a thread form that makes contact on the root and crest of the thread and the 3.degree. load bearing flank, but which permits a clearance on the trailing 10.degree. flank of as much as 0.007\" on smaller diameter tubulars. This clearance forms a leakage path that must be sealed by the solids in the thread lube. The API taper of 3/4\" per foot results in 1.984\" of imperfect (non-sealing threads) and, for a 7\", nominal, outside diameter casing corresponds to a buttress thread length of 4.200 inches.\nThe present invention represents an improvement on this API tubular connection and, as compared to that connection, provides a reduced maximum mated thread clearance, combines less imperfect threads with an increase in overall thread length to provide a greatly increased leakage path, and an advantageous taper mismatch which provides high bearing pressure and sealing within the made-up joint adjacent the narrow end of the male member where the coupling is strongest."}
-{"text": "Press fit sleeve bearings of both plastic and metal are well known. Metal sleeve bearings are typically pressed into a bore and then machined to the final internal dimension to accommodate the shaft that is intended to operate within the bearing. The post insertion machining is expensive and disrupts the normal flow of assembly of the product.\nThe plastic sleeve bearings have a diffenent drawback. In order for the sleeve to remain in the bore, there must be sufficient deformation in the bearing material to create the forces needed to retain the bearing in the bore. In plastics of the type used for bearings, the required deformation may be quite significant and therefore result in the deformation of the interior dimension. Since it is very diffficult to accurately predict the amount of internal deformation, due to variations in the tolerances of the bore hole and the exterior sleeve dimensions, it is preferable to form the bearing such that the bearing structure permits of controlled deformation. This then minimizes the requirement to perform post assembly operation on the bearing material. It is not very practical to perform a post insertion machining operation on the plastic bearing, since machining of plastic may not be precise and is is difficult to secure smooth bearing surfaces on plastic by machining. The molded skin of a plastic piece provides a superior bearing surface than that of a machined surface.\nPrior attempts to overcome the problem of using plastic materials as bearing sleeves have followed the approach of accommodating the compression of the bearing exterior without unduly effecting the interior bearing surface.\nU.S. Pat. No. 3,359,685 to L. M. Hodgen discloses a bearing sleeve which provides the bearing surface in a series of alternating lands with intermediate grooves, parallel to the axis of the bearing. The lands provide the bearing surface, while the grooves provide flexural relief to accommodate the compression of the material. The exterior of the sleeve is likewise a series of alternating lands and grooves. The lands of the exterior surface are radially aligned to the grooves of the interior surface, thereby permitting the distortion to be accommodated in the region of the groove of the interior surface with a minimum distortion of the bearing surface.\nU.S. Pat. Nos. 1,555,214 to C. W. Johnson and 3,515,417 to J. H. Bowman both use serrations or splines on the exterior surface of the bearing to accept the deformation caused by the insertion of the bearing into the bore of another member. In both of the patents, the splines or serrations extend the full length of the bearing.\nIn the case of the Johnson patent, the bearing is lined with a split sleeve which defines the interior dimension of the bearing. This sleeve may contract or expand. Dimensional stability is thereby sacrificed in the Johnson bearing.\nThe Bowman bearing appears more dimensionally stable on the interior but uses the flexure of the splines to effect a centering of the bearing surface.\nFor precision applications involving relatively light loads, plastic bearings are desirable since they may be molded to very tight tolerances. Cost is very reasonable, particularly if the molded dimensions may be used without further operations in the finally assembled product."}
-{"text": "1. Field of the Invention\nThe present invention relates to an intravascular prosthesis, and more particularly, concerns a percutaneously deliverable prosthesis suitable for intravascular reconstruction, reinforcement or repair.\n2. Description of the Prior Art\nBlood vessels of humans or animals undergo a natural degenerative process or are sometimes subject to weakness. In weakened blood vessels, aneurysms may occur. The degenerative effect on blood vessels may also cause a narrowing or constriction of the lumen of the vessel so that blood flow is restricted. In other degenerative situations, or for other reasons, clots or emboli may occur which, should they migrate within the intravascular system, could be very dangerous.\nWith respect to the aforementioned intravascular problems, surgical intervention has been the primary technique for providing relief. For example, aneurysm repair involves a surgical procedure in which an intraluminal vascular prosthesis is inserted into the damaged vessel to reconstruct the section that needs repair. For clogged blood vessels, the excision of thickened atheromatous areas of the vessel has been performed by an endarterectomy procedure. These and other intravascular therapy procedures of an invasive nature are not only risky, but are also costly.\nAngioplasty procedures, using expandable balloons, have been developed for widening the lumen of diseased, constricted blood vessels. Many of these angioplasty procedures are performed percutaneously so that the balloon is introduced into the blood vessel through a catheter inserted through the skin into the vascular system. After the inflation of the expandable balloon widens the clogged blood vessel, it is withdrawn from the blood vessel through the introducer catheter. Balloon catheters are also available in which the inflated balloon is detachable from the catheter once inflated within the blood vessel. The inflated, detached balloon occludes the blood vessel, and is therefore useful in such procedures as varicocele treatment. Such a detachable balloon catheter system, known as the MINIBALLOON.TM. catheter, is sold by Becton, Dickinson and Company, Paramus, N.J.\nAlthough inflatable and detachable balloon catheter procedures, performed percutaneously, are known and available for some intravascular therapy applications, surgery is still relied upon for other applications. Less invasive techniques are being sought for blood vessel repair reconstruction and filtering, as well as vessel occlusion, blood flow regulation or flow assist. The present invention is directed to a device which provides for minimal invasive methods of intravascular therapy. In particular, the present invention is directed to an intravascular prosthesis, percutaneously deliverable which is suitable for the reinforcement or reconstruction of weakened blood vessels, as well as aneurysm repair. Two copending patent applications, having a common assignee herewith, U.S. Ser. Nos. 772,217 now U.S. Pat. No. 4,705,517 and 772,218 both filed on Sept. 3, 1985, relate to those intravascular therapy applications involving blood filtering to prevent clot migration or emboli, and to blood vessel occlusion and blood flow regulation."}
-{"text": "Processes for removing light components such as hydrogen or nitrogen from a hydrocarbon gas mixture represent a problem in petroleum processing. Typically, the hydrocarbon gas mixture is a refining off-gas or a natural gas and the presence of the light components in the mixture limit further processing, or prevent sales of the gas mixture without the removal of the light components. Often these light components include hydrogen, nitrogen, carbon dioxide, carbon monoxide and mixtures thereof. A wide variety of hydrocarbon gas mixtures are found in petroleum refineries. Some streams are integral parts of specific processes. Often these streams are recycled from a fractionation column or a separator to a reactor. One such recycle stream may be an impure hydrogen stream which must be purified before returning the stream to the reactor. Other of these streams may be by-product streams of a major hydrocarbon conversion process such as a fluid catalytic cracking unit or an ethylene plant. Processes available for the rejection of the light components from these gas mixtures can employ solvent absorption, cryogenic separation, adsorption over molecular sieve adsorbents, or membrane separation. The choice of a suitable process depends upon many factors, some of which are the product purity that is desired, the gas recovery levels, available pressure losses, pretreatment requirements, off-gas composition, impact of reaction products remaining in the light component stream, and the turndown capability of the selected process.\nWhen natural gas is produced from a gas well or is recovered as associated gas from an oil well, the natural gas contains a number of light components which can reduce its heating value. Typically, a natural gas stream comprises nitrogen, methane, ethane, carbon dioxide, inerts, and C.sub.3.sup.+ hydrocarbons. In order to improve the heat content of the product natural gas, the concentration of nitrogen and other inerts must be reduced. This reduction of inerts is often required to meet the quality specifications of pipeline companies that transport natural gas from a well head or natural gas processing plant to the consumer. Typically the natural gas must meet the following specifications:\n______________________________________ Heat Content 900 to 1000 BTU Total Inerts (N.sub.2 + CO.sub.2) 7 mol % maximum Nitrogen 4 mol % maximum ______________________________________\nActual pipeline specifications vary somewhat depending upon the producer's contract for price and quality. In general, a specification for a higher heating value requires a reduced amount of light components such as nitrogen and carbon dioxide. Typically, natural gas at the well head can contain between 3 and 60 mol % nitrogen, between 0.1 and 10 mol % ethane, between 0.1 and 20 mol % C.sub.3.sup.+ hydrocarbons and between 0.1 and 20 mol % CO.sub.2 with the balance being methane. Light components, of which nitrogen is typically the major component, must be removed from natural gas to improve the heat content of the gas and to meet pipeline specification.\nCompositions of the raw gas and the amount of impurities that can be tolerated in the product generally determine the selection of the most suitable process for purification.\nCryogenic technology, which consists of several process variations, has traditionally been employed to separate nitrogen from natural gas. The operating principle of the process entails partially or fully liquefying the high nitrogen content feedstream under pressure and at very low temperatures (e.g., as low as -185.degree. C. (-300.degree. F.)). Afterwards, the partially or fully condensed stream is fractionated in one or two columns, which operate in tandem at two different pressures, to separate the feedstream into a rejected nitrogen stream overhead and a high methane content product stream at the bottom. After heat exchanges with the incoming feedstream, the rejected nitrogen stream is used for gas field operation, vented to atmosphere or reused. The bottom nitrogen-depleted product stream is regasified via process heat exchanges, recompressed, and delivered to the battery limits as an upgraded sales gas.\nU.S. Pat. No. 4,936,888 relates to the use of cryogenic technology in an integrated dual distillation system for rejecting nitrogen in concentrations as high as 75 mol % or more from gaseous methane in a distillation system employing a high pressure fractionator and a low pressure fractionator. The low pressure fractionator accepts a feed predominantly comprising liquid nitrogen under conditions to produce a high purity nitrogen overhead stream and a high purity methane bottom stream.\nU.S. Pat. No. 4,711,651 relates to a process for the separation of a high pressure gas stream such as refining gas in which the starting gas mixture is cooled and separated into a first vapor portion and a first liquid portion which is expanded to an intermediate pressure. The first vapor portion is further cooled and separated into a second vapor portion which may be further processed for ultimate recovery of a methane-rich product gas. Refrigeration is recovered from the mixed intermediate pressure stream.\nU.S. Pat. No. 5,051,120 relates to the cryogenic processing of a feed containing nitrogen and methane. The method is used in treating natural gas which is contaminated with nitrogen. An improved stripping column is provided in which the methane product is produced at a higher pressure than would otherwise be possible.\nOther process options have used non-cryogenic routes which employ solvent extraction. U.S. Pat. No. 4,832,718 relates to a process for processing a natural gas, a thermal or catalytic cracking gas or refining off-gas to produce a methane-rich product, a nitrogen-rich product, a hydrogen-rich stream or an olefin-rich product therefrom by solvent extraction. The process further relates to extractive flashing and extractive stripping using selective physical solvents. The process employs a solvent to separate N.sub.2 from hydrocarbons by extraction (i.e., absorption) at moderate temperatures (ambient to -30.degree. F.) and under pressure. As an N.sub.2 -rich gas stream and a solvent stream come in intimate contact in a packed column, the solvent dissolves nearly all of the hydrocarbons and a small amount of N.sub.2 in the feed gas stream, leaving behind a rejected gaseous N.sub.2 stream that contains the major fraction of the N.sub.2 previously present in the feed stream and a small amount of unrecoverable hydrocarbons. The dissolved hydrocarbons are subsequently released from the rich solvent by successive pressure reductions. The gases that evolve from the pressure reductions are recompressed and recombined into an upgraded sales gas stream. The regenerated lean solvent stream is recycled to the absorber. After refrigeration recoveries, the rejected N.sub.2 stream is available at a relatively high pressure and can be used for field applications, vented to the atmosphere, or reused.\nU.S. Pat. No. 4,623,371 relates to the recovery of hydrocarbons from natural gas containing acidic compounds and from 3-75 mol % nitrogen to provide up to three products: nitrogen-gas product, C.sub.1 -rich gas product, and a C.sub.2.sup.+ liquid product. The process extracts a natural gas stream with a physical solvent to produce a nitrogen stream and a methane-rich solvent stream. The methane-rich solvent stream is flashed to provide a stripped solvent stream and substantially all of the C.sub.1.sup.+ hydrocarbons as a stream of flashed-off-gases. The stripped solvent is recycled to the extraction step.\nU.S. Pat. No. 5,047,074 relates to a process for purging nitrogen from natural gas by passing the gas through an absorption column at elevated pressure and contacting the gas with an absorbent consisting primarily of a poly alpha olefin. The non-absorbed gas consisting predominantly of nitrogen is passed out of the absorption column to waste or to recovery. The resulting rich absorbent is desorbed to obtain the desorbed light hydrocarbons and the lean absorbent.\nU.S. Pat. 5,019,143 relates to the use of a particular group of solvents including paraffins, naphthenes, C.sub.6 -C.sub.10 aromatic compounds and dialkyl ethers of polyalkylene glycol to contact the gas feed mixture in a demethanizing absorber with a reboiler operating between 50 to 400 psig and a temperature of +10.degree. to -40.degree. F. for the separation of ethylene and lighter hydrocarbons from the feed with a distillation column to regenerate the rich solvent. U.S. Pat. Nos. 4,511,381 and 4,526,594 relate to the separation of natural gas from natural gas liquids with physical solvents. Natural gas liquids include hydrocarbons heavier than methane. The physical solvent used for the absorption step in U.S. Pat. No. 4,511,381 is described in terms of a solvent having a relative volatility of methane over ethane of at least 5.0 and a hydrocarbon loading capacity at least 0.25 standard cubic feet of ethane per gallon of solvent. The rich solvent stream is flashed and the gas fraction is compressed, cooled, and condensed to produce the natural gas liquids. U.S. Pat. No. 4,526,594 relates to the selective extraction of a stream of hydrocarbons that are heavier than methane from a natural gas stream with a physical solvent by selectively rejecting the consecutively lowest molecular weight portion of the extracted stream by the use of a second extraction step with a physical solvent to reject hydrocarbons heavier than methane, flashing the resulting rich solvent to separate the selected heavier hydrocarbons, and subsequently de-ethanizing the selected heavier hydrocarbons.\nU.S. Pat. No. 4,883,514 relates to a process for the removal of nitrogen from nitrogen-rich gases which contain more than 3 mol % nitrogen. The nitrogen-rich gas stream is contacted with a lean oil comprised of paraffins, aromatic or cyclohydrocarbons or mixtures thereof having molecular weight between 75 and 250 at temperature no lower than -40.degree. F. to produce a nitrogen stream as an overhead product and a bottoms methane-rich oil stream. The bottoms methane-rich oil stream is flashed to recover a methane-rich overhead gas product and a lean oil rich bottoms stream, and the lean oil stream is recycled to the absorption step. The patent teaches that the contacting or absorption may take place in a methane extraction column which includes a reboiler and is operable in an extractive stripping mode.\nIt is an object of this invention to provide a process for the rejection of nitrogen from natural gas streams using noncryogenic absorption with a hydrocarbon solvent. It is a further object of this invention to provide a process for the rejection of nitrogen from natural gas which can achieve gas recoveries of greater than 99 mol % and provides lower compression requirements, lower capital investment, and lower solvent losses than previous technology."}
-{"text": "The present invention relates in general to the field of patient supports, such as hospital beds. In particular, the invention relates to an improved control system that provides enhanced control of the patient support and accessories mounted at the patient support and, further, provides enhanced information related to the patient support and the patient and, further, allows for improved monitoring of the patient."}
-{"text": "Metal-insulator-metal (MIM) capacitors are often used in very-large-scale integration (VLSI) designs implemented on semiconductor chips, either as part of the primary circuit design or to provide decoupling capacitance for supply noise suppression. MIM capacitors are intended to improve capacitance density of the VLSI design while minimizing area overhead. FIG. 1 is a cross-section of an example semiconductor chip 100 that includes a MIM capacitor. The MIM capacitor comprises a capacitor top metal (CTM) node 104 and a capacitor bottom metal (CBM) node 106 formed on the top and bottom side, respectively, of a dielectric layer 108. The MIM capacitor layers reside between two metal layers MX and MX+1, on which die interconnection paths 110 are patterned. The MIM capacitor is insulated by an insulating layer 102. Conductive vias 112 are formed in the die to electrically connect circuit paths in the metal layers to the CTM node 104 or the CBM node 106 as needed.\nThe above-described description is merely intended to provide a contextual overview of current techniques and is not intended to be exhaustive."}
-{"text": "1. Field of the Invention\nThe present invention relates to a substrate cleaning apparatus and a substrate cleaning method that clean a substrate, in particular, a semiconductor wafer.\n2. Background of the Invention\nTypically, a resist on a semiconductor wafer substrate is removed (cleaned off) by immersing the semiconductor wafer in a chemical (resist remover) in a remover tank of a substrate cleaning apparatus.\nIn a conventional cleaning technique, for example, an SPM (sulfuric acid-hydrogen peroxide mixture) resist removing process, a mixture (SPM cleaning fluid) containing sulfuric acid and hydrogen peroxide in a certain proportion is first produced. Typically, to improve the resist removing capability, the SPM cleaning fluid serving as a remover is kept at a high temperature equal to or higher than 100 degrees Celsius, for example. Generally, hydrogen peroxide (H2O2) thermally decomposes into water (H2O) and oxygen (O2). In addition, hydrogen peroxide has a boiling point of 150 degrees Celsius and evaporates when heated to such a high temperature. Hydrogen peroxide tends to thermally decompose and evaporate at high temperatures in this way, so that the oxidation capability (resist removing capability) of the SPM cleaning fluid is lowered on the time passes.\nAccording to an existing measure against this, before introduction of a semiconductor wafer or the concentration of the SPM cleaning fluid is managed, a predetermined amount of hydrogen peroxide solution is additionally injected into an inner tank (wafer process tank) or an outer tank (overflow receiving tank) of a process tank at desired intervals, thereby maintaining the oxidation capability of the SPM cleaning fluid.\nA method of managing the hydrogen peroxide concentration uses a concentration monitor. However, since the concentration has to be always kept constant, the amount of usage of the hydrogen peroxide solution and sulfuric acid increases substantially.\nAccording to the method in which the hydrogen peroxide solution is introduced before introduction of the semiconductor wafer described above, it takes long before the hydrogen peroxide solution is sufficiently mixed with the circulating mixture to form a uniform SPM mixture containing the hydrogen peroxide solution and sulfuric acid.\nIn the case where the hydrogen peroxide solution is introduced into the inner tank as described above, the hydrogen peroxide solution having a low specific gravity cannot be effectively mixed with sulfuric acid and overflows into the outer tank, and thus, a desired resist removing capability cannot be obtained.\nSimilarly, in the case where the hydrogen peroxide solution is introduced into the outer tank, since the hydrogen peroxide solution is not effectively mixed, the hydrogen peroxide solution having a low specific gravity is not mixed with the sulfuric acid and remains separated in the upper part of the outer tank. Thus, it takes a predetermined time before a uniform mixture of sulfuric acid and the hydrogen peroxide solution is formed, or it is difficult to feed all the introduced hydrogen peroxide solution to a circulation piping because some of the hydrogen peroxide solution is discharged from the overflow pipe rather than being fed to the circulation piping.\nFurthermore, for example, a conventional substrate cleaning apparatus using the SPM cleaning fluid has a process tank that holds a mixture containing a hydrogen peroxide solution and sulfuric acid and is used for cleaning a substrate immersed in the mixture, circulation piping that extends between a primary side of the process tank on which the mixture is injected into the process tank and a secondary side of the process tank on which the mixture is discharged from the process tank and has a pump for causing circulation of the mixture, a filter disposed in the circulation piping for removing particles in the mixture, a chemical tank for additionally injecting a mixture of a hydrogen peroxide solution and sulfuric acid into the process tank, and a heater for heating the mixture in the chemical tank (see Japanese Patent Laid-Open No. 6-342780, for example).\nHowever, in the conventional substrate cleaning apparatus, the mixture is heated by the heater to a temperature higher than a temperature required in the process tank before the mixture is injected into the process tank. Thus, the SPM cleaning fluid heated to the high temperature thermally decomposes and evaporates, so that the hydrogen peroxide concentration of the SPM cleaning fluid decreases. As a result, a desired SPM cleaning fluid cannot be obtained in the process tank.\nAs described above, the conventional technique described above has a problem that the SPM cleaning fluid cannot have a required resist removing capability, and the amount of the hydrogen peroxide solution additionally injected into the SPM cleaning fluid cannot be optimized."}
-{"text": "The advantages of combinatorial chemistry for the rapid generation of chemical compounds for pharmaceutical screening are well established. The single greatest strength of combinatorial chemistry is that it makes possible the generation of libraries of huge numbers of compounds in a relatively short period of time.\nSeveral methods have been employed for the preparation of such libraries, e.g., solution-phase synthesis; solid-phase synthesis on polymer beads or other divided supports; synthesis on soluble, precipitatable polymer supports; and synthesis on planar supports. All of these methods are capable of generating large libraries (i.e., libraries containing a large number of compounds), but none of them are amenable to rapid screening of the libraries for binding activity, biological activity, or other desirable properties.\nCurrent methods for screening libraries include those in which individual library members, or small groups of members, are assayed in microtiter plates, e.g., by screening for a desired activity or for binding to a specific binding partner, such as a receptor or antibody or other ligand, but the number of compounds that can be assayed at once is on the order of 102 to 104: one compound per well in a 96-well microtiter plate screens at most 96 compounds, while twenty compounds per well in a 384-well plate screens about 7680 compounds. The latter approach, while permitting higher throughput, requires secondary screening to identify the active species in any given well. For example, screening a library containing all of the 3.2 million possible pentapeptides which could be made from the twenty natural amino acids (i.e., a 3.2-million-member library) by these methods would require 500 to 3400 plates. Screening a library of the 64 million possible hexapeptides would require 10,000 to 68,000 plates. Robotic systems are available for microtiter plate assays, but screening a large library by such a method remains a massive undertaking. Furthermore, in a pharmaceutical discovery environment this represents only one of the many assays an organization might wish to conduct.\nAn alternative to testing individual compounds is the testing of complex mixtures, with various \u201cdeconvolution\u201d strategies being employed to deduce the active species. These strategies have in common the re-synthesis and re-testing of successively less-complex mixtures. In addition to the great effort involved, the testing of complex mixtures is limited by the low concentration of any individual species in the mixture, and is susceptible to false positive results from the additive effects of large numbers of weakly active species. In practice, deconvolution has had limited success in probing large libraries of compounds. See L. Wilson-Lingaro, J. Med. Chem., 39, 2720-2726 (1996); and D. A. M. Konings, J. Med. Chem., 39, 2710-2719 (1996), and references therein, for a discussion of deconvolution strategies.\nIdeally, the probing of a combinatorial library would be conducted in a single operation, with the active members of the library being in some way \u201cpointed out\u201d of the vast population of compounds by the assay. Two of the current methods meet this requirement. Both methods employ solid-phase synthesis, and both require that the library members remain attached to the solid support. In one method, a library of compounds bound to polystyrene beads is prepared by the split-mix method. The library is assayed in a single batch, by being exposed to a molecule of interest, such as a receptor, enzyme, or other specific binding partner. Any beads to which the molecule binds are visualized (e.g., by a colorimetric assay), and beads so identified are selected and the structure of the library member attached to the bead is determined. This can be done by sequencing if the compound is a peptide or nucleic acid, as described, e.g., in U.S. Pat. No. 5,382,513 (incorporated herein by reference). The structure of the compound on the bead may in some cases be deduced from spectroscopic evidence (see, e.g., U.S. Pat. No. 5,382,513), or by decoding a chemical tag that reveals the chemical history of the bead, as described in patent application WO 95/24186 (incorporated herein by reference). The method is in principle capable of screening very large libraries, limited only by the number of beads one is willing to examine. In practice, libraries of 104 to 106 members can be dealt with in this fashion.\nThe second approach involves physically locating a compound or compounds in a spatially addressable array of compounds on a planar support. In this approach, a compound's identity is revealed by its location in the array. One method of this type employs an array of compounds generated by light-directed synthesis, as first disclosed by Fodor et al. in Science, 251, 767-773 (1991), in which a fraction of sites on a planar support carrying photo-detachable protecting groups is exposed to light through a photolithographic mask, and the fraction of sites thus deprotected are functionalized with a specific monomer or building block, itself carrying a photo-detachable protecting group. The process is repeated with the mask in a different position or orientation, or with a different mask, and a second monomer or building block is attached to the support and/or to the first monomer residues. After numerous such cycles, with careful attention to the pattern of masking, an array of compounds is built up on the support. The final array is completely deprotected, and exposed to the ligand of interest. Binding of the ligand is visualized by immunofluorescence, using antibodies against the ligand which are tagged with a fluorescent dye. Under a fluorescence microscope, any location in the array to which the ligand has bound is visible as a fluorescent area, and the x-y coordinates of the area reveals the identity of the library member to which the ligand was bound. This technology, as applied to polypeptide and oligonuceotide synthesis, is known as Very Large Scale Immobilized Polymer Synthesis, or VLSIPS. It is described in U.S. Pat. Nos. 5,143,854, 5,413,939, 5,424,186, and 5,527,681, all of which are incorporated herein by reference.\nThe photolithographic method of synthesis, however, is cumbersome, and requires a substantial investment in very specialized equipment. A further investment in a fluorescence microscope or a specialized scanner is required for the assay, and highly skilled technicians are required, at least to conduct the synthesis aspect of the process. The scale of library synthesis is limited by the size of the masks and by the translational reach of the scanning device, which together limit the accessible surface for synthesis to a few square centimeters. In practice this technique is presently limited to arrays of 104 to 105 compounds. For these reasons the method is not routinely employed; see G. Jung and A. G. Beck Sickinger in Angewandte Chemie, 31, 367 (1992).\nAs an alternative to photolithography, the use of directed laser light to conduct light-directed synthesis has been described in U.S. Pat. Nos. 4,719,615 and 5,318,679 (both of which are incorporated herein by reference). In the latter patent, a rectangular array support is either held stationary or translated, while a laser beam is scanned across the array by means of a rotating mirror, in the manner of an ordinary laser printer. This provides an alternative to the photolithographic masking approach, but fluorescence microscopy is still relied upon as an assay method.\nAnother alternative method for the synthesis of spatially addressable arrays utilizes ink-jet printing technology to spray micro-droplets of reagent solutions onto a substrate. This method is disclosed in U.S. Pat. Nos. 5,474,796 and 5,449,754, both of which are incorporated herein by reference. This method is in theory capable of preparing arrays of 107 compounds, but a method of indexing such large arrays is not disclosed. Again, fluorescence microscopy is the preferred means of conducting an assay.\nThere remains a need for reliable methods of generating very large, very high-density arrays of chemical compounds, on the order of about 108 or more compounds, along with a method of rapidly screening such arrays for chemical properties of interest, such as binding to antibodies, cellular receptors or other ligands, catalytic activity, or inhibition of enzymes."}
-{"text": "The present invention relates to connectors containing fusible materials to assist in forming a connection and more particularly to such connectors which, during the heating of the fusible material, form part of a circuit, the temperature of which is autoregulated at about the Curie temperature of magnetic material included in the circuit at least during the heating operations.\nU.S. Pat. No. 3,243,211 discloses a connector containing a fusible material so that upon insertion of an object to be joined to the connector or insertion into the connector of two members to be joined, and upon heating of the connector, the fusible material is caused to melt and contact said object or objects and upon cooling, effect a bond. The connector may also include a heat recoverable member whereby the liquified fusible material is bonded and caused to contact the object or objects while in the fluid state. This device requires an external heat source to melt the fusible material, such as hot air or an infrared radiant source.\nThe difficulty with the device of the patent is in the danger of overheating the objects to be soldered or otherwise bonded as well as adjacent objects. In the electronics art, for instance, overheating of delicate integrated circuits is a problem as is overheating of circuit boards, mastics, resins, heat shrinkable polymers, glues, potting compounds, all of which can be destroyed by the application of excessive heat. Further, the device has little utility for joining wires, tubes or members which are large effective heat sinks, since the large amount of heat required cannot be readily transferred through the heat shrinkable sleeve without damaging it."}
-{"text": "As computer systems have become faster in processing speed, improvement in the reliability of such systems has become an issue.\nThis means that a technique is needed that can make systems operate stably without shutting them down even when a failure has occurred.\nMulti-CPU systems, which are equipped with plural CPUs, can operate stably.\nThe techniques described below are proposed as techniques for the detection of failures in multi-CPU systems.\nPatent Document 1 discloses a technique by which operations are monitored in units of processors. In the technique of Patent Document 1, data is transmitted between processors that are connected via a communications device, and one processor monitors responses from the other processors in order to detect failures.\nPatent Document 2 discloses a technique by which an abnormal operation in a system is monitored in units of instructions so that the system can recover from the failure. Patent Document 2 discloses a method in which the occurrence of a failure in a processor is monitored by the processor itself and the processor performs a resetting operation using an automatic reset occurrence circuit and a failure element holding circuit when the system can recover by resetting.\nThe methods disclosed by the above described patent documents have the problems described below.\nIn the method of Patent Document 1, processors are monitored by communications between processors themselves. This configuration causes a time gap of several minutes between the occurrence of a failure and the recognition of the failure.\nIn the method of Patent Document 1, information used for investigating the cause is collected after the recognition of a failure, and this means that the information collected for the investigation is not exactly from the moment at which the failure occurred. This makes it difficult to use the information for investigating the cause, and thereby the investigation tends to take a long time period or sometimes the investigation itself is prevented.\nThe time scale of monitoring a system (in units of minutes or seconds) and the time scale of a CPU (in units of nanoseconds) are significantly different from each other, resulting in difficulty in understanding the status of the system when a failure occurred, and thereby a long time is required to investigate the cause.\nAlthough there is a technique of collecting information on changes in the status of a CPU in the form of log information, effective information sometimes cannot be collected using this technique because of the limitation on the volume of log information.\nAccording to the technique of Patent Document 2, the monitoring and recovering is performed in units of instructions, and the system is reset in order to restart when the restarting of the firmware can make the system recover from a state involving failure.\nIn the technique of Patent Document 2, the monitoring and recovering is performed in units of instructions, making it difficult to detect a failure in a program or the like in an OS (Operating System)."}
-{"text": "A) Field of the Invention\nThis invention relates to an electronic musical apparatus, and more in detail, an electronic musical apparatus that can display lyrics and a chord name on other electronic musical apparatus.\nB) Description of the Related Art\nIn an electronic musical apparatus that has an automatic musical performing function such as an electronic musical instrument, when music data including lyrics data is reproduced, it is well-known that an external displaying apparatus displays lyrics via a video-out device (image data output circuit), for example refer to JP-A 2002-258838.\nIn the above-described prior art, lyrics corresponding to music data are output to an external apparatus as image data, and lyrics can be displayed on a separated displaying device and a displaying device that has a large screen.\nIn the prior art, however, image data (image signals) for displaying lyrics is generated based on lyrics data and image data is transmitted to an external apparatus via a video-out device, this kind of apparatus is expensive since the video-out device is generally expensive."}
-{"text": "While just about every computer user owns their own printer and is capable of producing high quality documents, the ability to produce such documents in high volume and with special finishing features, such as binding, is still within the purview of the commercial print shops and corporate copy departments. High volume, finished production of documents is typically referred to as production printing. A production printer is a printing device capable of rapid production of large volumes of documents. Typically these printers have high paper handling capacity, the ability to draw on multiple media types from multiple sources and the ability to automatically finish a document such as by adding a binding. Despite the automation provided by the production printer and the proliferation of computer technology, especially in the area of desktop publishing, production printing is still a complicated and often manual process.\nIn a typical print shop, customers bring in original documents which they want turned into a finished product such as a bound booklet, a tri-fold brochure or a tabbed three ring bound notebook. In addition, they typically need a large volume of the finished product, for example, one thousand brochures. The combination of the original documents plus the instructions for producing the finished product is called a \u201cjob\u201d. The documents can be brought in either in hard copy or electronic form, such as on floppy disk, compact disc or tape or can be transmitted to the print shop over a network such as the Internet.\nAfter handing over the documents to the clerk, the customer relays his instructions for preparing the finished product. The clerk will note these instructions on a \u201cticket\u201d or \u201cjob ticket\u201d. The job ticket is typically a piece of paper with all of the instructions written on it for producing the finished product. As mentioned above, this is known as job. The job will then be handed to an operator, who runs the production printer, to produce the finished output. The operator's job is to prepare the document for production, load the appropriate materials, such as paper stock and binding materials, into the production printer and ensure that the finished output is correct.\nWhile the job of the operator seems simple, there are many issues which quickly complicate it. Often, the documents provided by a customer are not ready to be run on the production printer. Some documents provided by a customer are merely raw manuscripts requiring basic formatting, such as margins, typography, etc. Other documents may be formatted but such formatting might not take into account the requested binding. For example, the text of the document is too close to the margin, therefore, when the finished product is bound, some of the text will be obscured. Some documents, such as books, require special care so that, for example, the first page of every chapter appears on the front of a page, also known as imposition. Other forms of imposition include booklet/pamphlet imposition or n-up imposition. Or the customer may bring in multiple documents and ask that these \u201cchapters\u201d be assembled into a book, with a cover and binding.\nOther issues which complicate the production printing job are determining and loading the correct media into the production printer. Often, jobs will require many different paper types, such as different stock weights or different colors. In addition, some jobs require the insertion of tab stock at specific points within the document. Still other jobs may require the adding of a Bates number or other annotation to the document.\nWith such a complicated production process to produce finished output, errors are bound to occur, such as loading the wrong paper stock in the printer or setting a margin too close to a binding. Production printers run at very high speeds, often producing output greater than 1 page per second therefore, errors in the finished output may not be caught before a significant amount of time and resources have been wasted.\nAccordingly, there is a need for an efficient system and method for managing the production printing workflow."}
-{"text": "Golf is enjoyed by a wide variety of players\u2014players of different genders and dramatically different ages and/or skill levels. Golf is somewhat unique in the sporting world in that such diverse collections of players can play together in golf events, even in direct competition with one another (e.g., using handicapped scoring, different tee boxes, in team formats, etc.), and still enjoy the golf outing or competition. These factors, together with the increased availability of golf programming on television (e.g., golf tournaments, golf news, golf history, and/or other golf programming) and the rise of well known golf superstars, at least in part, have increased golf's popularity in recent years, both in the United States and across the world.\nGolfers at all skill levels seek to improve their performance, lower their golf scores, and reach that next performance \u201clevel.\u201d Manufacturers of all types of golf equipment have responded to these demands, and in recent years, the industry has witnessed dramatic changes and improvements in golf equipment. For example, a wide range of different golf ball models now are available, with balls designed to complement specific swing speeds and/or other player characteristics or preferences, e.g., with some balls designed to fly farther and/or straighter; some designed to provide higher or flatter trajectories; some designed to provide more spin, control, and/or feel (particularly around the greens); some designed for faster or slower swing speeds; etc. A host of swing and/or teaching aids also are available on the market that promise to help lower one's golf scores.\nBeing the sole instrument that sets a golf ball in motion during play, golf clubs also have been the subject of much technological research and advancement in recent years. For example, the market has seen dramatic changes and improvements in putter designs, golf club head designs, shafts, and grips in recent years. Additionally, other technological advancements have been made in an effort to better match the various elements and/or characteristics of the golf club and characteristics of a golf ball to a particular user's swing features or characteristics (e.g., club fitting technology, ball launch angle measurement technology, ball spin rates, etc.). Further technological advancement in golf club design has also involved the incorporation of various types of monitoring devices or sensors in the golf club. Many such designs, however, have been cumbersome and unreliable. In addition, further processing of the data recorded by the sensors has been limited or not performed in a suitable manner to be most useful to golfers.\nWhile the industry has witnessed dramatic changes and improvements to golf equipment in recent years, there is room in the art for further advances in golf club technology. Thus, while golf equipment according to the prior art provide a number of advantageous features, they nevertheless have certain limitations. The present invention seeks to overcome certain of these limitations and other drawbacks of the prior art, and to provide new features not heretofore available."}
-{"text": "A gate-array type integrated circuit is a semiconductor device which has basic circuit elements arranged in a pattern on a chip (the \"master slice\"). The circuit elements can be interconnected during manufacture to achieve a desired functionality. Thus, a gate array can be customized.\nCircuit elements typically found in a gate array include storage elements, such as latches, and logic elements, such as AND gates. A \"path\" is defined as a signal transmission passageway between and including circuit elements. There are various types of paths. For example, there are closed paths between storage elements; open paths from a storage element to an output terminal or from an input terminal to a storage element; and through paths from an input terminal to an output terminal. The \"input\" and \"output\" terminals are external terminals (or \"leads\") of the semiconductor circuit device.\nFIG. 11 illustrates sample paths (1) through (4). For example, path (1) includes latches 901 and 902 and gates 909 and 910. Wire portions 917, 918 and 919 connect these elements as illustrated. Each element and wire has an inherent capacitance, which causes a propagation delay as the signal moves from latch 901 to latch 902. These delays are illustrated in the bar graphs of FIGS. 15 and 16. The portion of the propagation delay attributable to each particular element or wire portion is designated in FIGS. 15 and 16 by boxes bearing the reference numeral of the particular element or wire portion.\nAs shown in FIG. 11, latches 901 and 902 receive timing signals T.sub.0 and T.sub.1, respectively. It is essential that the signal propagate from latch 901 to latch 902 during the interval between timing signals (the \"delay reference value\"). Otherwise, the signal will not have reached latch 902 by the time latch 902 is strobbed by timing signal T.sub.1. The delay reference value is shown in FIGS. 15 and 16 by a broken line.\nAs can be seen in FIG. 15, some paths, such as path (4), have propagation delays which exceed the delay reference value. Others, such as path (3), have propagation delays which are less than the delay reference value. As is known, the current between elements can be varied to adjust the propagation delay attributable to the capacitive effects of the wire portions. Specifically, the greater the current, the shorter the propagation delay.\nTo control flow of current between elements, each element includes at its output an emitter-follower circuit, comprised of a plurality of resistors and transistors. As explained in Japanese Post-Examination Application No. 64-4340, these resistors and transistors are selectively coupled to obtain an emitter-follower circuit having the desired current characteristic. By increasing current, the propagation delays can be shortened, so that the aggregated delays for a given path do not exceed the delay reference value.\nThe conventional method of varying an emitter-follower current in accordance with the capacitance of a load circuit will be described with reference to FIGS. 11, 15 and 16. As explained above, FIG. 15 shows the delays in the paths (1)-(4) of FIG. 11 before emitter-follower current is adjusted. The delays illustrated in FIG. 15 include delays in elements and delays on signal lines. FIG. 16 shows the delays after the emitter-follower current has been adjusted in accordance with conventional techniques.\nFor example, consider paths (1) and (2) of FIG. 11. Both paths (1) and (2) have the same number of elements. However, signal line 921 of path (2) is longer than signal lines 917-919 of path (1). Effective capacitance (that is, delay) caused by line 921 is larger than that caused by lines 917-919. Thus, the overall delay of path (2) is greater than the delay of path (1), as shown by the bar graphs of FIG. 15. By increasing the emitter-follower current of a logic element 911, the delay can be shortened, as shown in FIG. 16.\nIn the conventional method, such an adjustment is conducted to all long signal lines, which sometimes causes an unnecessary increase of power consumption. For example, (3) of FIG. 11 includes a long line 923. However, because path (3) has few elements, its overall delay is less than the delay reference time, as shown in FIG. 3. Nevertheless, in accordance with conventional practice, the emitter-follower current of storage element 905 is increased, thus reducing the overall delay of path (3) further below the delay reference time (as shown in FIG. 16), and resulting in additional unnecessary power consumption.\nPath (4) of FIG. 11 includes three elements 914-916 and three short lines 925-928. Because signal lines 925, 926 and 928 are short, their effective capacitances are small. Increasing the emitter-follower current of elements 915 cannot sufficiently reduce the delay of line 927, and it is impracticable to reduce delay on short lines 925, 926 and 928. Therefore, even after adjustment, the delay of path (4) exceeds the delay reference value.\nThus, the conventional technique simply adjusts the emitter-follower current of each element in accordance with the capacitance of the load circuit associated with that element. When the capacitance of the load circuit is large, the element's emitter-follower current is increased, even through the element's path may already have a sufficiently short delay. On the other hand, when the capacitance of the load circuit is small, the current value is not altered, even though the element may be part of a path which has an aggregate delay that exceeds the delay reference value.\nAs described above, the conventional technique reduces the delay between the elements having large capacitive loads. Therefore, the delay of certain individual elements may be reduced. The conventional technique is not effective, however, for reducing the delay in paths which have a larger number of elements connected by short signal lines. Therefore, the speed of the whole semiconductor circuit is limited by the path having the longest delay, and thus limiting the overall speed of the integrated circuit.\nFurther, the conventional method does not reduce the delays attributable to the elements, themselves, as opposed to the delay attributable to the lines between elements. Further, the conventional method does not minimize power consumption."}
-{"text": "Memory devices are typically provided as internal storage areas in a computer. The term memory identifies data storage that comes in the form of integrated circuit chips. There are several different types of memory used in modern electronics, one common type is RAM (random-access memory). RAM is characteristically found in use as main memory in a computer environment. RAM refers to read and write memory; that is, you can both write data into RAM and read data from RAM. This is in contrast to ROM (read-only memory), which permits you only to read data. Most RAM is volatile, which means that it requires a steady flow of electricity to maintain its contents.\nComputers almost always contain a small amount of ROM that holds instructions for starting up the computer, typically called a basic input output system (BIOS). Unlike RAM, ROM generally cannot be written to by a user. An EEPROM (electrically erasable programmable read-only memory) is a special type of non-volatile ROM that can be erased and programmed by exposing it to an electrical charge. EEPROM comprise a large number of memory cells having electrically isolated gates (floating gates). Data is stored in the memory cells in the form of charge on the floating gates. Charge is transported to or removed from the floating gates by specialized programming and erase operations.\nYet another type of non-volatile memory is a Flash memory. A Flash memory is a type of EEPROM that can be erased and reprogrammed in blocks instead of one byte at a time. A typical Flash memory comprises a memory array, which includes a large number of memory cells. Each of the memory cells includes a floating gate field-effect transistor capable of holding a charge. The data in a cell is determined by the presence or absence of the charge in the floating gate. The cells are usually grouped into sections called \u201cerase blocks.\u201d The memory cells of a Flash memory array are typically arranged into a \u201cNOR\u201d architecture (each cell directly coupled to a bitline) or a \u201cNAND\u201d architecture (cells coupled into \u201cstrings\u201d of cells, such that each cell is coupled indirectly to a bitline and requires activating the other cells of the string for access, but allowing for a higher cell density). Each of the cells within an erase block can be electrically programmed in a random basis by charging the floating gate. The charge can be removed from the floating gate by a block erase operation, wherein all floating gate memory cells in the erase block are erased in a single operation. It is noted that other types of non-volatile memory exist which include, but not limited to, Polymer Memory, Ferroelectric Random Access Memory (FeRAM), Ovionics Unified Memory (OUM), Magnetoresistive Random Access Memory (MRAM), Molecular Memory, Nitride Read Only Memory (NROM), and Carbon Nanotube Memory.\nEach erase block of a Flash memory device contains a series of physical pages that are typically each written to a single row of the Flash memory array and include one or more user data areas and associated control or overhead data areas. The control/overhead data areas contain overhead information for operation of physical row page and the user data area each overhead data space is associated with. Such overhead information typically includes, but is not limited to, erase block management (EBM) data, sector status information, or an error correction code (ECC). ECC's allow the Flash memory and/or an associated memory controller to detect data errors in the user data area and attempt to recover the user data if possible.\nMany of the internal operations of volatile and non-volatile memories require that the memory perform data comparisons. Typically this data comparison is performed in the context of comparing data that has been read from the memory array with the data that was expected to be read in order to find discrepancies. The internal operations requiring data comparison include, but are not limited to, data write, block erasure, and memory testing. A problem with many modern memory devices and arrays is that, because of their increasing storage density levels and increasing size of each physical page/row of data of the array, this data comparison within the memory device has become a time consuming task and can affect the speed of operation, data throughput, and testing time required.\nFor the reasons stated above, and for other reasons stated below which will become apparent to those skilled in the art upon reading and understanding the present specification, there is a need in the art for circuits and methods to allow fast and efficient comparison of large amounts of data within memory devices."}
-{"text": "The present invention relates to a picture frame for framing pictures, posters, prints, signs and the like. More particularly, it relates to such a picture frame for the inexpensive, yet aesthetically pleasing and practical framing of a wide variety of pictures for both residential and commercial use.\nIn my earlier patent for a picture frame, U.S. Pat. No. 4,669,209, I have described a picture frame which is very inexpensive and so easy to assemble that anyone can do so with aesthetic and precise results. The picture frame is designed to eliminate the need for a picture glass or other transparent covering and semi-rigid backing, normally used for most picture frames. It was specifically intended for framing posters, which tend to be relatively large, and prints. Artistic decor utilizing posters and large prints has come into vogue in recent years. Where such poster-type pictures are inexpensive, it is often times not worthwhile to provide the same with a glass covering and a semi-rigid backing, together with a frame, for the entire assembly since these various extra elements often cost far in excess of the original poster to be framed. Thus, on many occasions, one sees attractive poster-like pictures which are tacked or otherwise fastened to a wall for the purpose of exhibiting the picture and maintaining it flat. Such securement of a poster-like picture to a wall not only causes damage to the wall but very often tends to be unsightly, in that tacks or tape, or other objects are utilized in fastening the picture to the wall.\nTo overcome this problem, my earlier patent provided a picture frame for poster-like pictures which included at least a top picture edge support, and preferably a bottom picture edge support, and fastening elements which fasten the respective edges of the poster-like picture frame to the top and bottom supports. Each picture edge support includes a flat back plate which extends substantially the length of the edge of the picture being framed and a substantially C-shaped portion or member, one leg of which is connected to the outer edge of the back plate, while the other leg extends downwardly toward the flat back plate to define a gap between the terminus thereof and the back plate. The picture edge support is preferably formed of extruded plastic. The fastening elements are also generally C-shaped and resiliently engaged around the generally C-shaped member of the edge support with the inside leg extending into the gap beyond the terminus of the inside leg of the C-shaped member and the flat back plate, so as to press the edge of the poster-like picture inserted into the gap against the back plate of the edge support. The invention also contemplates the provision of edge supports for each side edge of the poster-like picture, wherein the fastening elements are in the form of elbow-shaped corner pieces, which engage the adjacent ends of adjacent edge supports so as to secure or fasten the poster-like picture to the flat back plates of the supports at the corners of the picture. While quite satisfactory in use, I have found that for certain applications a more rigid or stiffer picture frame is needed. In addition, the picture frame is not suitable for joining multiple pictures or posters together. It is also not particularly suitable for double-sided signs, such as may be often used in store windows and the like, where it is desired that the customer can read the same either on the outside of the store or on the inside thereof."}
-{"text": "1. Field of the Invention\nEmbodiments of this invention are directed to the field of intravascular medical devices, and more particularly to the field of catheters such as angioplasty, neurological and guide catheters, among others, which may be used in various medical procedures such as percutaneous transluminal angioplasty (PTA), percutaneous transluminal coronary angioplasty (PTCA) as well as in procedures involving the placement of medicines and medical devices within the body.\nSome embodiments of the invention are directed to all forms of catheters which may be advanced through a body lumen or vessel. Some examples of catheters are over-the-wire (OTW) catheters, such as are described in U.S. Pat. No. 5,047,045; single-operator-exchange (SOE) balloon catheters, such as are described in U.S. Pat. Nos. 5,156,594 and 5,549,552. Other examples of catheters which may incorporate the unique features of the present invention may include rapid-exchange catheters, guide catheters, etc.\n2. Description of the Related Art\nIntravascular diseases are commonly treated by relatively non-invasive techniques such as PTA and PTCA. These angioplasty techniques typically involve the use of a balloon catheter. In these procedures, a balloon catheter is advanced through the vasculature of a patient such that the balloon is positioned proximate a restriction in a diseased vessel. The balloon is then inflated and the restriction in the vessel is opened. In other uses a catheter may be used to delivery an endoprosthesis such as a stent, graft, stent-graft, vena cava filter or other implantable device or devices herein after collectively referred to as a stent or stents. Where a stent is to be delivered into a body lumen the catheter may include one or more inflatable portions or balloons. Typically, the stent is retained in the predelivery state about the catheter shaft, or a portion thereof such as a balloon, by crimping and/or through the use of a retaining mechanism such as sleeve, sheath or sock.\nBalloons and balloon catheters may be particularly useful for the delivery of expandable, implantable medical devices such as stents, grafts, stent-grafts, vena cava filters, hereinafter referred to cumulatively as stents. Stents and catheters used in their delivery are commonly used and as such their structure and function are well known.\nA stent is a generally cylindrical prosthesis introduced via a catheter into a lumen of a body vessel in a configuration having a generally reduced diameter and then expanded. In its expanded configuration, the stent supports and reinforces the vessel walls while maintaining the vessel in an open, unobstructed condition.\nIn order to properly position a stent and/or balloon within a body lumen, the catheter must be advanced through the narrow confines of the body. Typically the balloon and/or stent is located near the distal end of the catheter. In order to advance the distal end of most prior catheters further in to a body lumen, the inner shaft or catheter member is utilized to transmit force to the distal end. However, the inner is typically soft and flexible which often results in poor transmission of push.\nThe present invention, in accordance with the various embodiments presented herein, addresses the shortcoming of poor push transmission common to many catheters.\nWithout limiting the scope of the invention, a brief summary of various embodiments of the invention is set forth below. Additional details of the summarized embodiments of the invention and/or additional embodiments of the invention may be found in the Detailed Description of the Invention below.\nA brief abstract of the technical disclosure in the specification is provided as well for the purposes of complying with 37 C.F.R. 1.72.\nThe entire content of all of patents or other references listed within the present patent application are incorporated herein by reference."}
-{"text": "Field of the Invention\nOne or more embodiments of the invention are related to the field of containers. More particularly, but not by way of limitation, one or more embodiments of the invention relate to a cup lid with integrated container that enables simultaneous or intermittent access of the contents of the container and attached cup without disengagement of the cup lid from the cup. Additionally, an independent drop-in container may reside within the cup lid cavity such that after partially consuming the contents of the independent drop-in container it may be resealed with an additional lid and removed from the lid cavity and transported to another location such as a car or home.\nDescription of the Related Art\nStandard cup lids are simple covers that do not include an integrated container. Rather, known lids cover the contents of a cup which forms a closed container in combination with the cup itself. Known containers that couple with cups include food containers that fit onto the top of yogurt cups for example. Known containers have to be removed from the yogurt cup and then flipped over and opened before the contents of the container and cup may be accessed. It is generally not possible to access the contents of the cup while also accessing the contents of the container without first disengaging the container from the cup. Additionally, food containers that attach to yogurt cups in an upside-down position have a limited food-volume capacity. In such cases, as the yogurt example shows, the food-container walls narrow as they proceed upward toward the bottom of the upside down container. Other known devices having a container or shelf combined with a lid have limitations which makes these devices impractical to use. One category of devices includes a container combined with a cup, but utilizes a hole in the middle of the lid. This makes it impossible to store relatively circular items, i.e., non-ring or non-annular items having no central hole, in the container, such as hamburgers, cookies or muffins for example. Another category of device includes a container combined with a lid, but does not allow for simultaneous access of the contents of the cup and the container at the same time, and does not allow for the container to be resealed or a drop-in container to be removed from the container. Other devices that include drop-in functionality require removal of the container before accessing the contents of the cup. Yet these devices do not contemplate a drop-in container that is configured to fit into the arm rest of a movie theater seat. Other devices have relatively small peel containers for pills such as mints and are not suitable for larger food items. Another category of devices utilizes dividers in the cup with access on each side of the cup. None of the known devices enable a container to be disengaged from the lid of the cup while retaining the lower lid on the cup. No known devices have a non-permanent or male/female bottom oriented coupling system for coupling a container with the lid. Furthermore, there are no known rotational covers that enable or disable access to the liquid and/or solid in the cup as desired by the user.\nKnown containers that couple with bottles include gift containers that fit onto the top of bottles for example. It is generally not possible to access the contents of the bottles while also accessing the contents of the gift containers without disengaging the gift container from the bottle and then disengaging the lid of the bottle.\nThus simultaneous or intermittent access of the contents of known cups or bottles and of the contents of an attached container is not possible. This makes for difficult drinking/eating of coffee, soda, snacks, popcorn, etc., in malls, fast food restaurants, theaters, amusement parks, sports stadiums or in any other venue. For example, this makes it difficult to eat and drink food in a theater or stadium with one cup-holder per seat.\nFor at least the limitations described above there is a need for a cup lid with integrated container."}
-{"text": "1. Field of the Invention\nThis invention relates to fluid dispensing systems which automatically dispense fluid to animals only during predetermined periods of time.\n2. Description of the Prior Art\nExperiments and controlled feeding often require maintaining animals on a fluid restriction regimen. These restrictions may require limiting the amount of fluid an animal may consume, the time period during which the animal may consume the fluid, or some combination of both of these restrictions. Controlled feeding programs in which a fluid is available to an animal at the same time every day present scheduling difficulties. Typically, glass bottles containing fluid are manually placed in the animal's cage each time access to the fluid is required. This manual procedure often produces cumbersome work schedules for experimentors and laboratory personnel, particularly when the experiments last several weeks or months, including weekends and holidays.\nU.S. Pat. No. 3,294,066 describes an animal feeding system wherein a time actuated valve causes a predetermined amount of water to enter a feeding unit where it will mix with a dry food material and the weight of such mixture will downwardly displace the feeding unit so as to expose a feeding nipple portion, thus allowing the animal to feed from the unit. This system, unlike the present invention, lacks a provision for draining or retrieving any fluid remaining when an animal feeding interval is over and it requires that a disposable feeding bag be replaced after each use. Additionally, this system differs from the present invention in that it is designed to dispense a fixed amount of fluid and is not adaptable to experiments where the watering time interval, rather than the quantity of fluid to be consumed, is to be restricted.\nU.S. Pat. No. 3,720,185 discloses an automatic animal feeding system in which a level sensing device controls the amount of liquid delivered to a chamber where it is mixed with a solid food material and then dispensed at individual feeding points. This system utilizes automatic valves to dispense the food but is designed to provide a predetermined amount of food to an animal rather than to allow the animal access for a limited period of time."}
-{"text": "The present invention relates to a beverage preparation machine, and particularly though not exclusively to an espresso-type coffee machine that includes an outlet nozzle for steam or other liquids such as milk. These nozzles, called steam nozzles, are immersed into the beverage. They subsequently require a cleaning of at least their outside surface that has come in contact with the beverage. More particularly, the invention relates to a beverage preparation machine comprising:                a body;        a movable steam nozzle connected to the body by a moving device;        a beverage preparation area in which a cup collecting the beverage can be placed;        a cleaning area which is separate from the preparation area and in which is disposed at least one rinsing container for the nozzle, the moving device being capable of moving the nozzle between the preparation and cleaning areas.        \nCleaning a steam nozzle, for example after preparing a cappuccino or a latte, is time-consuming since conventionally, it requires a wet cloth or a container of water and cleaning product. Consequently, after the preparation of a beverage, the steam nozzle is quite often left covered with residue, particularly in the case of a home machine used by several people. Cleaning it becomes even more difficult after traces of milk have dried on the outside wall of the nozzle. It is then practically essential to let the nozzle soak in a container containing a cleaning solution, typically water mixed with a cleaning product like a detergent. Yet the presence of a container, a source of hot water if the machine does not have one, and particularly a bottle of cleaning product near the machine is not very satisfactory. Aside from the space issues, it is necessary to make sure that the bottle of cleaning product is actually present, that it still contains some liquid, that this liquid is measured out correctly by the users, and that it is not within reach of children. To avoid altering the flavor or appearance of a beverage prepared afterwards, it is also important to thoroughly rinse or wipe off the nozzle after it is cleaned.\nAll these inconveniences mean that the cleaning of a steam nozzle is often neglected, even though in certain machines it is relatively easy to manually move the steam nozzle in order to immerse it into a receptacle placed next to the cup rest.\nMoreover, there are systems, known particularly from the document WO03091152A, for automatically cleaning a nozzle similar to a steam nozzle. These systems place the nozzle in a relatively closed chamber and dispense a cleaning solution under pressure through the nozzle. The vigorous agitation of the liquid around the nozzle cleans it. However, the main purpose of these systems is to clean out the internal conduits through which a product, like milk, that is susceptible to bacterial contamination has circulated. These systems require solenoid valves and complex and expensive dosing devices. Aside from having to switch the nozzle between a beverage supply and a supply of water under pressure, it is necessary to provide for the delivery into the conduits of a cleaning product in precise doses, on the order of a few percent of the water flow. These systems operating under pressure are therefore costly, especially since it is necessary to guarantee high degrees of watertightness and operational reliability in order to avoid dispensing a residue of cleaning solution along with the beverage."}
-{"text": "1. Field of the Invention\nThe present invention relates to a composition of polybutene-1 and a process for preparing the same and, more particularly, to a composition and a process for the preparation of polybutene-1 which has less gel content and less fish eyes and which can preferably be used in the field of film.\n2. Description of Related Art\nA homopolymer of polybutene-1 or a copolymer of butene-1 with another olefin, as generally called polybutene-1, has been used for pipes for warm water due to its excellent creep resistance at high temperatures or in the film field due to its high orientation at a low degree of stretching. Such conventional polybutene-1, however, has the drawback that it is remarkably poor in moldability or forming performance because of its low speed of crystallization. Hence, various attempts have so far been made to solve the drawback of conventional polybutene-1.\nIn order to accelerate the cystallization speed, for example, U.S. Pat. No. 4,321,334 proposes using a nucleating agent in preparing polybutene-1. In this process, however, it is so difficult to disperse the nucleating agent in polybutene-1, that it cannot be said to produce the effect to a sufficient degree. Particularly, when polybutene-1 obtained by this process is formed to film, the film suffers from the disadvantage that fish eyes are caused to thereby give a poor appearance.\nFurther, Japanese Patent Unexamined Publication (kokai) No. 123,607/1980 discloses a process for preparing crystalline polybutene-1 having a high bulk density by subjecting a small amount of .alpha.-olefin other than butene-1 to preliminary polymerization in the process for preparing polybutene-1. As this process is developed with the object to provide polybutene-1 with high bulk density by using a small amount of the other .alpha.-olefin for preliminary polymerization, the above-identified patent publication is silent about the problems of gel content and fish eye. Turning now to specific examples of the description of the patent publication, it can be noted that they do not use a molecular weight modifier, such as hydrogen or the like, during preliminary polymerization. Hence, it can be presumed that a polymer obtainable by this preliminary polymerization may have a too high molecular weight, thereby causing gel and fish eye in large numbers."}
-{"text": "Mobile cellular networks are typically engineered by operators to provide diverse sets of voice and data services to groups of users sharing network resources. In a standard configuration, each of the networked radio coverage cells (also known as \u201csectors\u201d) are arranged to provide static contiguous wireless coverage within a regional service area. Users of a mobile cellular network connect to the network via user equipment terminals (UEs), which attach wirelessly to one or more of the cells.\nThe network and the UEs support handover services, which provide continuous and uninterrupted user-perceived communication as users move between coverage cells. For example, the network may handover a UE from a first cell to a second cell when the UE moves from the first cell into the second cell.\nBecause of a random distribution of users and a varying demand for voice and data services, traffic loads vary from cell to cell. Adjoining or overlapping cells may have dramatically different loads over time.\nWhen a cell experiences more demand for resources than it can satisfy, the cell is \u201ccongested.\u201d When the cell is congested, overloaded shared resources within the cell may be allocated on a fair basis among the UEs attached to the cell. However, each UE is allocated fewer resources when the cell is congested than when the cell is not congested. Those reduced allocations during congestion may be insufficient for certain high data throughput services, such as video services or timely large file delivery services, e.g., downloading email attachments. Congestion is therefore a condition to be avoided whenever possible.\nAccordingly, mobile networks, when possible, \u201cload balance\u201d cells by handing over UEs between serving cells. Load balancing prevents scenarios, for example, where one cell is operating in congestion, and another neighboring or overlapping cell is idle. Load balancing tends to maximize the usage of the aggregate capacity of the cellular network, improve the economic use of deployed infrastructure, and improve quality of experience for users of the mobile network.\nFor example, operators may trigger load balancing between cells whenever a numerical threshold difference exists in the aggregate traffic loads between adjoining or overlapping coverage cells that have one or more operating UEs in a common coverage area. For example, by forcing or biasing handover of UEs from a higher traffic load cell to a lower traffic load cell, the aggregate traffic load levels between the two cells are driven towards equilibrium, and general congestion may be lessened.\nExisting methods, however, currently lack novel methods of triggering and executing load balancing to more efficiently and equably reduce overall network congestion that may be aligned with Self-Organizing-Network (SON) principles."}
-{"text": "Collapsible canopy frames often include a plurality of telescoping legs, each having one or more X-shaped scissor assemblies extending therebetween. A canopy covering, such as a cloth or leather covering, is disposed above, and supported by, the collapsible canopy frame. The X-shaped scissor assemblies are moveable relative to the telescoping legs to adjust the collapsible canopy frame between an expanded position and a collapsed position.\nIn the expanded position, the collapsible canopy frame provides a temporary shelter. In the collapsed position, the collapsible canopy frame can be more readily transported. Typically, collapsible canopy frames are transported by placing the collapsible canopy frame on a separate wheeled structure, such as a wheeled platform. However, collapsible canopy frames are often relatively heavy and it is therefore desirable to minimize any upward lifting that is required during transportation of a collapsible canopy. Accordingly, a need exists for a wheel assembly and/or components related thereto for attachment to a collapsible canopy frame that facilitates transportation of the collapsible canopy frame."}
-{"text": "The present invention relates to an apparatus for automatically examining and inspecting white resin powder of, for instance, vinyl chloride resins, ABS resins and MBS resins for the presence of foreign substances.\nWhen preparing resin powder or forwarding the powdery product from a manufacturing plant, the product is in general subjected to quality inspection for various predetermined properties. The results are sent back to the manufacturing plant for the improvement of production processes or they are used in the denoration of the quality in order to afford convenience to the destination or the consignee.\nIn case of, for instance, vinyl chloride resin powder, one of the items for quality inspection thereof is to inspect the resin powder for the presence of foreign substances as contaminants. Conventionally, the inspection of the powdery resin product for the presence of foreign substances has been carried out by spreading a constant amount of the vinyl chloride resin powder and searching for the presence of pigmented foreign substances (such as those colored black, brown and/or red) thus to count the number of foreign substances present therein. However, the results widely vary depending on the visual power and the nature of each inspector. Accordingly, this method never provides objectively measured values, requires much labor and time and accordingly, has low efficiency.\nIncidentally, there has recently been used an industrial apparatus for automatic visual inspection which makes use of a video camera. An example of such apparatuses is one which is commercially available from Hajime Sangyo Co., Ltd. under the trade name of FF 4000. This apparatus serves to detect, for instance, defects, stains and foreign substances on the image of a powdery resin sample photographed by the video camera and to specify the positions of the foreign substances on the photograph through blinking modulated bright spots on a monitor. This inspection apparatus permits automatic inspection of a subject for the presence of, for instance, foreign substances and can provide measured values free of any scattering irrespective of inspectors. Moreover, the apparatus does not require much labor and time and ensures efficient inspection.\nHowever, when the aforementioned industrial apparatus for automatic visual inspection which makes use of a video camera is used in the inspection of resin powder easily electrostatically charged such as vinyl chloride resin powder for the presence of colored foreign substances as contaminants thereof, it is often observed that the resin powder to be inspected is not uniformly dropped on a belt conveyor or is unevenly distributed on the belt conveyor due to electrostatic charges generated on the resin powder per se. For this reason, the surface of the spread resin powder layer becomes severely uneven and this leads to a large scattering in the measured values. Consequently, the apparatus does not provide any reproducible result."}
-{"text": "Computer systems are designed to process a variety of applications, each comprised of software instructions. With the increasing complexity of applications, software instructions have become longer, thus requiring an increased amount of time for the software instructions to be executed.\nSoftware instructions are executed by a microprocessor which is the key working unit of a computer system. The methods that have been used to increase speed in the personal computer have generally centered on maximizing the efficiency with which a single microprocessor can process instructions. Limits are being reached on single microprocessor processing speed. To address this constraint, multiple microprocessors have been combined to operate in parallel within computer systems. Such multiple microprocessor systems, such as symmetrical multiprocessor (i.e., SMP) systems, allocate processing tasks among the multiple parallel microprocessors.\nUpon powering on of an SMP system, each microprocessor is initialized, and a primary microprocessor is selected from among the microprocessors to take charge of bringing up the SMP system. The selection of the primary microprocessor can be done by a software algorithm, or by a hardware locking mechanism. Depending on the number of microprocessors in an SMP system, the time required for initialization increases with the number of microprocessors in the SMP system. The selection of the primary microprocessor can vary upon every time the SMP system is powered on."}
-{"text": "The utility of aqueous polymeric dispersions (also called aqueous polymeric emulsions) in the preparation of paints, coatings, adhesives and caulks or sealants is well known. Water based or water borne polymer dispersions are often preferred because of theft relatively low cost, ease of application and relatively low amounts of volatile organic compounds (VOC) contained therein. There is an increasing need for higher solids dispersions which will provide faster setting times for use on high speed production equipment. High solids adhesive bases may also find use as replacements for conventional hot melt (100% solids) material which requires elevated temperatures with consequent expenditures of considerable energy. In addition to the need for high solids dispersion for such applications, it is also essential that the dispersions remain sufficiently low in viscosity that they can be applied using conventional equipment. Considerable effort has been expended to provide aqueous dispersions of polymeric materials which are characterized by a high solids content.\nPressure-sensitive adhesives find use in a wide variety of applications, such as automotive, aerospace, construction and electrical markets, either in the form of tapes or as adhesive coatings on other backings. As the society becomes more aware of the significance of environmental protection, industrial products such as adhesive that is harmful to the ecological environment are gradually being replaced and eliminated.\nCoated surfaces, especially those coated with modern low VOC coatings containing 100 g/L. VOC or less, such as moldings and panels often become visibly discolored after contact with water repeatedly or over extended periods of time. This problem is enhanced when high polymer content coatings are used on the surfaces. Other issues that can occur with painted surfaces are blistering and surfactant leaching, wherein water-soluble components are extracted from the coatings and deposited on the coated surface. Minimizing the amount of water-soluble ingredients has been used to reduce surfactant leaching. However, the water-soluble components are generally substituted with volatile organic compounds (VOCs), resulting in an environmentally undesirable product.\nSurfactants have widely been used as emulsification, dispersion, cleaning, wetting and foaming. Emulsifiers for emulsion polymerization, which are used upon producing polymers by emulsion polymerization, are known not only to take part in polymerization-initiating reactions and polymer-forming reactions but also to affect the mechanical stability, chemical stability, freezing stability, storage stability and the like of the resulting emulsions. Further, they are also known to give significant effects on physical properties of the emulsions, such as particle size, viscosity and foaming potential and, when formed into films, physical properties of the films, such as waterproofness, weatherability, adhesion and heat resistance.\nSurfactants can also be used as reactants, often called reactive surfactants. Surfactants used in polymerization are also called polymerizable and/or copolymerizable surfactants. Among the reactive surfactants, those containing one or more phenyl ether groups as hydrophobic groups have found wide-spread utility for their excellent properties such as emulsifying property, dispersing property, and polymerization-stabilizing property. Thus, surfactants that are based on alkyl phenol ethoxylates (APES) have been widely used in emulsion polymerization. In recent years, however, a concern has arisen about a potential problem that nonyl phenol ethoxylates may show false hormone effects on organisms to disrupt the endocrine system, that is, the so-called endocrine problem has arisen, so that research has also been conducted in efforts to provide replacements for the reactive surfactants containing one or more phenyl ether groups.\nAccordingly, there is a need in the art for an aqueous dispersion and/or emulsion with ultra-high solids but environmentally friendly in preparations and applications."}
-{"text": "This invention relates to non-lethal weapons; and more particularly, to a electrical stun gun and electrically conductive liquids."}
-{"text": "Modern bathroom design favors a look with planer surfaces and configurations which minimize the visibility of unattractive space. Nowhere is this more noticeable than in the design of bathtubs and tub accessories. Traditional tubs had a shape matching their name, showing a curved exterior surface and exposed piping. Modern design provides more boxlike shape which hides piping and unattractive space beneath and around the bottom of the tub. A common way to provide this shape as part of a new tub system or often over the top of an older tub is to install a tub skirt. A tub skirt provides a flat, aesthetically pleasing front surface.\nA tub skirt, while hiding unattractive space, also prevents access to this space. As a remedy, tub skirts have been provided with access openings. These openings provide access for cleaning or tub repair and to access motors and plumbing on Jacuzzi type tubs. When such activities are not occurring, the openings are covered. Common accessories used to cover these openings are decorative panels. U.S. Pat. No. 5,940,906 to Halloran discloses a skirt frame with a detachable panel. The panel is held in place with Velcro tabs. U.S. Pat. No. 5,804,898 to Kapp et al. discloses a skirt frame with a detachable panel. Velcro is used for attachment. U.S. Pat. No. 5,208,924 to Smith et al. discloses a skirt frame and mechanically attachable panel. These panels are attached by screws with caps.\nSkirts are often provided without panels. In other cases, the panels provided with the skirts are not aesthetically pleasing and retrofitting is desired by the consumer. Common ways to attach aftermarket or retrofit panels in order to match new bathroom colors or the like is with velcro pads because they can be attached with an adhesive. However, velcro tends to lose holding force over time and cannot be cleaned easily. This lack of holding force is especially noticeable in Jacuzzi type tubs which include motors which vibrate the tub skirt slightly. Panels which are held in place with screws and caps are difficult to remove. Removal is a time consuming process. What is desired is a simple system which remains hidden from view, but can be used repeatedly without a loss in effectiveness."}
-{"text": "1. Field of the Invention\nThe present invention relates to an integrated circuit, and more particularly to a high voltage integrated circuit.\n2. Description of the Prior Art\nA variety of power supplies and motor drivers utilize the bridge circuits to control a power source to the load. The bridge circuit normally has a high-side transistor connected to the power source and a low-side transistor connected to ground. A common node between the high-side transistor and low-side transistor is coupled to the load. As transistors are controlled to alternately conduct, the voltage of the common node swings between the voltage of the power source and the ground. Therefore the control of high-side transistor requires a charge pump circuit and/or a floating driver in order to fully turn on the high-side transistor. In recent development, many floating circuits are disclosed in U.S. Pat. No. 6,344,959 (Milazzo), U.S. Pat. No. 6,781,422 (Yang) and U.S. Pat. No. 6,836,173 (Yang).\nFIG. 1 shows a high-side transistor drive circuit, in which a circuit 10 is the floating driver. A capacitor 40 provides a supply voltage to the floating driver 10. A voltage VD charges the capacitor 40 through a diode 45 once the low-side transistor 30 is switched on. The ground reference of the capacitor 40 is pulled to the level of the voltage source VIN when the high-side transistor 20 is turned on. This happens, because turning on the high-side transistor 20 shifts the bridge circuit into a low impedance state."}
-{"text": "The present invention relates to a salt and associated hydrates of racemic 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid, processes for their preparation, pharmaceutical compositions comprising them, and their use in antibacterial therapy.\nEP 688772 (corresponding to Korean Patent Laid open Publication No 96-874) discloses novel quinoline(naphthyridine)carboxylic acid derivatives, including anhydrous 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid of formula I, having antibacterial activity. \nAccording to the invention there is provided 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate.\n7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate (hereinafter referred to as xe2x80x98the methanesulfonatexe2x80x99) may be obtained as an anhydrate or a hydrate (i.e., methanesulfonate.nH2O).\nHydrates of the methanesulfonate wherein n is in the range of from 1 to 4 are preferred. Particular hydrates of the methanesulfonate, which may be mentioned, are those in which n is 1, 1.5, 2, 2.5, 3, 3.5, and 4. Particularly preferred compounds are those in which n is 1.5 or 3, with n=1.5 being most preferred.\nThe moisture content of the methanesulfonate hydrates varies with the hydration number (n) of the hydrated molecule. The methanesulfonate has a molecular weight of 485.5. Thus the calculated moisture content of hydrates where n is 1, 1.5, 2, 2.5, 3, 3.5, and 4 is 3.6%, 5.0%, 6.9%, 8.5%, 10.0%, 11.5%, and 12.9%, respectively. However, the actual moisture content of the methanesulfonate hydrates may differ from the calculated value depending on various factors including recrystallization conditions and drying conditions. The observed moisture content for the methanesulfonate hydrates where n is 1, 1.5, 2, 2.5, 3, 3.5, and 4 is shown in Table 1:\nIt is possible to mix methanesulfonate hydrates having different moisture contents together to give materials having intermediate moisture contents.\nPreferred methanesulfonate hydrates have a moisture content of from 4 to 6% or from 9 to 11%, especially a moisture content of from 4 to 6%.\nThe methanesulfonate has been observed to exist as a stable hydrate over a range of hydration numbers (n). Stability of the hydrate refers to its resistance to loss or gain of water molecules contained in the compound. The methanesulfonate hydrates maintain a constant moisture content over an extended relative humidity range. The n=3 hydrate has a constant moisture content at a relative humidity of from at least 23 to 75%, and the n=1.5 hydrate has a constant moisture content at a relative humidity of from 23 to 64% (see FIGS. 3 and 4). In contrast, moisture absorption by the anhydrate varies greatly with relative humidity.\nBoth the methanesulfonate anhydrate and n=3 hydrate undergo transition to the n=1.5 hydrate in aqueous suspension indicating that the latter is thermodynamically more stable. The n=1.5 hydrate is a sesquihydrate at 11 to 64% of relative humidity. Above 75% relative humidity, it takes up water over 10% and its X-ray diffraction pattern changes. The hydrate (another form of n3 having different physicochemical properties from the n=3 hydrate of Example 2) obtained from n1.5 hydrate at 93% relative humidity is not stable at lower relative humidity, and it converts back to n=1.5 hydrate at a relative humidity below 75%.\nSince the moisture content of the anhydrate changes readily depending on the environment (e.g., relative humidity, formulation additives, etc.), it may require careful handling during storage or formulation, with operations, such as quantifying procedures, being performed in a dry room. The hydrates do not change in moisture content easily and hence products, which are stable under prolonged storage and formulation may be obtained. The hydrate can be tableted without the addition of a binder since the water contained in the compound itself acts as a binder, whereas it may not be possible to tablet the anhydrate at a similar pressure.\nThe present invention also provides a process for the preparation of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate and hydrates thereof which comprises reacting 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid with methanesulfonic acid and crystallizing the resulting methanesulfonate from solution, and where desired or necessary adjusting the hydration of the compound.\nThe methanesulfonate and its hydrates may be prepared by the addition of methanesulfonic acid to the free base which may be prepared as described in EP 688772. Preferably, 0.95 to 1.5 molar equivalents of methanesulfonic acid are added to the free base, or 1 molar equivalent of methanesulfonic acid dissolved in a suitable solvent is added to the free base. Suitable solvents for the preparation of the methanesulfonate and its hydrates include any solvent in which the methanesulfonate is substantially insoluble, and the suitable solvents include C1-C4 haloalkanes, C1-C8 alcohols and water, or mixtures thereof. Dichloromethane, chloroform, 1,2-dichloroethane, methanol, ethanol, propanol and water, or mixtures thereof, are preferred solvents. If necessary, the free base may be heated in the solvent to facilitate solution before methanesulfonic acid is added, and alternatively the methanesulfonic acid may be added to a suspension, or partial suspension, of the free base in the solvent. Following addition of the methanesulfonic acid, the reaction mixture is preferably allowed to stand or is stirred for 1 to 24 hours at a temperature of from about xe2x88x9210 to 40xc2x0 C. The resulting methanesulfonate is obtained as a solid, which can be isolated by filtration or by removal of the solvent under reduced pressure.\nDifferent hydrates may be obtained by altering the recrystallization conditions used in the preparation of the methanesulfonate, and such conditions may be ascertained by conventional methods known to those skilled in the art.\nThe present invention also provides a process for the preparation of a hydrate of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate comprising exposing the methanesulfonate anhydrate or a solvate thereof to a high relative humidity.\nThe methanesulfonate anhydrate or solvate thereof is preferably exposed to a relative humidity of at least 75%.\nThe methanesulfonate anhydrate or solvate thereof may be exposed to high relative humidity by passing humidified nitrogen gas through the methanesulfonate anhydrate or solvate thereof or by standing the methanesulfonate anhydrate or solvates thereof under a high relative humidity.\nThe humidified nitrogen gas used in this process, for example nitrogen gas having a humidity of at least 75%, may be made by conventional methods. In this process it is desirable to maintain the temperature in the range above which moisture condensation could occur. Also, particularly in large scale production, it is preferable to stir the sample thoroughly while the humidified nitrogen gas is passed through. When the hydrate is prepared by standing the methanesulfonate anhydrate or solvate thereof under a high relative humidity, for example a relative humidity of at least 75%, it is preferable to spread the sample as thinly as possible in order to raise the conversion efficiency.\nThe solvates of methanesulfonate anhydrate which may be used in the process according to this aspect of the present invention include solvates with one or more organic solvents. Preferred solvents include C1-C4 haloalkanes and C1-C8 alcohols, for example those selected from the group consisting of ethanol, dichloromethane, isopropanol and 2-methyl-2-propanol.\nSolvates of the methanesulfonate anhydrate are novel. Thus according to a further aspect of the invention there is provided a solvate of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate with one or more organic solvents.\nThe solvates of the methanesulfonate are prepared by recrystallization and controlled by the condition of recrystallizing system.\nThe methanesulfonate and its hydrates exhibit the same potent antibacterial activity as the corresponding free base disclosed in EP 688772. The methanesulfonate and its hydrates also exhibit desirable physicochemical properties including improved solubility and constant moisture content regardless of the ambient relative humidity when compared to the free base and other salts thereof. The methanesulfonate and its hydrates thus exhibit greater ease of handling, quality control and formulation than the free base and other salts thereof.\nAs mentioned above the methanesulfonate and its hydrates exhibit antibacterial activity. The methanesulfonate and its hydrates may be formulated for administration in any convenient way for use in human or veterinary medicine, according to techniques and procedures per se known in the art with reference to other antibiotics, and the invention therefore includes within its scope a pharmaceutical composition comprising 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate or a hydrate thereof together with a pharmaceutically acceptable carrier or excipient.\nCompositions comprising the methanesulfonate or hydrate thereof as active ingredient may be formulated for administration by any suitable route, such as oral, parenteral, or topical application. The compositions may be in the form of tablets, capsules, powders, granules, lozenges, creams, or liquid preparations, such as oral or sterile parenteral solutions or suspensions. Tablets and capsules for oral administration may be in unit dose presentation form and may contain conventional excipients, such as binding agents, for example, hydroxypropyl methyl cellulose, hydroxypropyl cellulose, syrup acacia, gelatin, sorbitol, tragacanth, or polyvinylpyrrolidone; fillers, for example, microcrystalline cellulose, lactose, sugar, maize-starch, calcium phosphate, sorbitol, or glycine; tableting lubricants, for example, magnesium stearate, talc, polyethylene glycol, or silica; disintegrants, for example, sodium starch glycolate, cross-linked polyvinylpyrrolidone, or potato starch; or acceptable wetting agents such as sodium lauryl sulfate. The tablets may be coated according to methods well known in normal pharmaceutical practice. Oral liquid preparations may be in the form of, for example, aqueous or oily suspensions, solutions, emulsions, syrups or elixirs, or may be presented as a dry product for reconstitution with water or other suitable vehicle before use. Such liquid preparations may contain conventional additives such as suspending agents, for example, sorbitol, methyl cellulose, glucose syrup, gelatin, hydroxyethyl cellulose, carboxymethyl cellulose, aluminum stearate gel or hydrogenated edible fats; emulsifying agents, for example, lecithin, sorbitan monooleate, or acacia; non-aqueous vehicles (which may include edible oils), for example, almond oil, oily esters, glycerine, propylene glycol, or ethyl alcohol; preservatives, for example, methyl or propyl p-hydroxybenzoate or sorbic acid; and, if desired, conventional flavoring or coloring agents. Suppositories will contain conventional suppository base, e.g., cocoa-butter or other glyceride.\nFor parenteral administration, fluid unit dosage forms are prepared utilizing the compound and a sterile vehicle, water being preferred. The methanesulfonate or hydrate thereof, can be either suspended or dissolved in the vehicle, depending on the vehicle and concentration used. In preparing solutions the methanesulfonate or hydrate thereof can be dissolved in water for injection and filter sterilized before filling into a suitable vial or ampoule and sealing. Advantageously, agents such as local anaesthetic, preservative and buffering agents can be dissolved in the vehicle. To enhance the stability, the composition can be lyophilized and the dry lyophilized powder sealed in a vial, and an accompanying vial of water for injection may be supplied to reconstitute the powder prior to use. Parenteral suspensions are prepared in substantially the same manner except that the methanesulfonate or hydrate thereof is suspended in the vehicle instead of being dissolved and sterilization cannot be accomplished by filtration. The methanesulfonate or hydrate thereof can be sterilized by exposure to ethylene oxide before suspending it in the sterile vehicle. Advantageously, a surfactant or wetting agent is included in the composition to facilitate uniform distribution of the methanesulfonate or hydrate thereof.\nThe methanesulfonate or hydrate thereof may also be formulated as an intramammary composition for veterinary use.\nThe composition may contain from 0.1% to 100% by weight, preferably from 10 to 99.5% by weight, more preferably from 50 to 99.5% by weight of the active ingredient measured as the free base, depending on the method of administration. Where the compositions comprise dosage units, each unit will preferably contain from 50-1500 mg of the active ingredient measured as the free base. The dosage as employed for adult human treatment will preferably range from 100 mg to 12 g per day for an average adult patient (body weight 70 kg), for instance 1500 mg per day, depending on the route and frequency of administration. Such dosages correspond to approximately 1.5 to 170 mg/kg per day. Suitably the dosage is from 1 to 6 g per day.\nThe daily dosage is suitably given by administering the active ingredient once or several times in a 24-hour period (e.g., the active ingredient up to 400 mg may be administered once a day). In practice, the dosage and frequency of administration which will be most suitable for an individual patient will vary with the age, weight, and response of the patients, and there will be occasions when the physician will choose a higher or lower dosage and a different frequency of administration. Such dosage regimens are within the scope of this invention.\nThe present invention also includes a method of treating bacterial infections in humans and animals comprising administering a therapeutically effective amount of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate or a hydrate thereof.\nIn a further aspect, the present invention also provides the use of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate or a hydrate thereof for the manufacture of a medicament for treating bacterial infection.\nThe methanesulfonate and its hydrates are active against a broad range of Gram-positive and Gram-negative bacteria, and may be used to treat a wide range of bacterial infections including those in immunocompromised patients.\nAmongst many other uses, the methanesulfonate and its hydrates are of value in the treatment of skin, soft tissue, respiratory tract and urinary tract infections, and sexually transmitted diseases in humans. The methanesulfonate and its hydrates may also be used in the treatment of bacterial infections in animals, such as mastitis in cattle."}
-{"text": "Convolution is commonly performed on signals in many contexts, including the fields of sound, still image, video, lithography, and radio (radar) signal processing. Typically, the signals to be convolved are pattern signals. Each of the expressions \u201cpattern\u201d and \u201cpattern signal\u201d is used herein in a broad sense to denote a one-dimensional sequence or two-dimensional (or higher dimensional) array of data words (which can be, but need not be pixels). Typically, the data words comprise binary bits, and the convolution is performed in discrete fashion on the binary bits using software, digital signal processing circuitry, custom hardware, or FPGA systems (field programmable gate array based computing systems).\nThe term \u201cdata\u201d herein denotes one or more signals indicative of data, and the expression \u201cdata word\u201d herein denotes one or more signals indicative of a data word.\nThe motivations for implementing convolution rapidly, even when processing data indicative of very large patterns, are myriad. The present invention was motivated by the need for proximity correction in the field of lithography. In such problems, one attempts a two-dimensional convolution between data indicative of a large pattern \u201cp\u201d (where the pattern is a pixel array) and a diffusion kernel \u201cd\u201d. Often the kernel \u201cd\u201d is a Gaussian or a superposition of Gaussians, or is otherwise a smooth kernel. More specifically, the present invention grew out of attempts to establish a suitable \u201cO(NN\u2032)\u201d algorithm (an algorithm requiring not more than on the order of NN\u2032 multiplications and additions) for convolving a two-dimensional pattern comprising NN\u2032 pixels, where each of N and N\u2032 is very large) with a Gaussian kernel (or other smooth kernel) such that the convolution is exact or very close to exact.\nThe objective in performing proximity correction (in the field of lithography) is to generate a \u201craw\u201d optical signal (or \u201craw\u201d electron beam signal) which can be input to a set of reflective or refractive optics (or electron beam optics), in order to cause the output of the optics to produce a desired pattern on a mask or wafer. To determine the characteristics of a raw optical signal (or raw electron beam signal) that are needed to produce the desired pattern on the mask or wafer, a deconvolution operation is typically performed on a very large array of pixels (which determine a pattern \u201cp\u201d) in order to correct for the well known proximity problem. In the case of electron beam lithography, the proximity problem results from electron scattering in the substrate (mask or wafer) being written. Such scattering exposes broadened areas on the substrate to electrons (i.e., an area surrounding each pixel to be written in addition to the pixel itself), with the scattering effectively broadening the electron beam beyond the beam diameter with which the beam is incident on the substrate.\nIn nearly all proximity correction schemes, such a deconvolution operation includes at least one convolution step. Accordingly, in performing typical proximity correction, a very large array of pixels (determining a pattern \u201cp\u201d) must be convolved with a diffusion kernel. Although such a convolution is typically performed on a pattern comprising a very large array of binary pixels, this restriction is not essential in the following discussion and is not essential to implementation of the invention. Indeed, the invention can implement convolution on data indicative of any pattern \u201cp\u201d with a smooth convolution kernel \u201cd\u201d having characteristics to be described below.\nFor data indicative of a pattern \u201cp\u201d and a convolution kernel \u201cd\u201d we consider the cyclic convolution:                     (                  d          \u2062                      \u00d7            C                    \u2062          p                )            n        =                  \u2211                              i            +            j                    \u2261                      n            \u2062                          (                              mod                \u2062                                  \u2003                                \u2062                N                            )                                          \u2062                        d          i                \u2062                  p          j                      ,where \u00d7C denotes that the convolution operator has cyclic character, and an acyclic convolution which differs only in the indicial constraint and range:                     (                  d          \u2062                      \u00d7            A                    \u2062          p                )            n        =                  \u2211                              i            +            j                    =          n                    \u2062                        d          i                \u2062                  p          j                      ,where \u201c\u00d7A\u201d denotes that convolution operator has acyclic character.\nFor simplicity, we restrict much of the discussion herein to one-dimensional cases (in which the pattern p is an ordered set of N data values and the kernel is an ordered set of M values). Despite this, it should be appreciated that in typical embodiments of the invention, the pattern is two-dimensional (a two-dimensional array pjk of data values determines the pattern) and the summation defining the convolution (a summation which corresponds to either one of the summations set forth in the previous paragraph) is over index k as well as index j of the array pjk. In the case of a two-dimensional pattern p determined by an N by N\u2032 array of data values, the indices n, i, j and domain lengths N in the formulae set forth in the previous paragraph are 2-vectors.\nIn one-dimensional cases, the result of the cyclic convolution has length N (it comprises N data values), and the result of the acyclic convolution has length M+N\u22121.\nIt is standard that a cyclic convolution d\u00d7Cp can be cast as an equivalent matrix-vector product:d\u00d7Cp\u2261Dp,where D is the circulant matrix of d (hereinafter the \u201ccirculant\u201d of d), whose 1-dimensional form is defined (assuming that N is greater than 3) as:   D  =            (                                                  d              0                                                          d              1                                                          d              2                                                          d              3                                            \u22ef                                              d                              N                -                1                                                                                        d                              N                -                1                                                                        d              0                                                          d              1                                                          d              2                                            \u22ef                                              d                              N                -                2                                                                          \u22ee                                              \u2003                                                          \u2003                                                          \u2003                                                          \u2003                                            \u22ee                                                              d              1                                                          d              2                                                          d              3                                                          d              4                                            \u22ef                                              d              0                                          )        .  Therefore, conventional methods for cyclic convolution can be cast in the language of matrix algebra. Acyclic convolution can be obtained with similar matrix manipulations.\nFor the sake of simplicity, we will use the symbol \u00d7 hereinbelow to denote convolution having either acyclic or cyclic character. In most of the discussion the symbol will refer to convolution having acyclic character. Those of ordinary skill in the art will recognize that, given a specified acyclic convolution, a corresponding cyclic convolution can be implemented by slight modification of the parameters (e.g., the boundary conditions and definition of the circulant of the kernel) that determine the acyclic convolution.\nAbove-referenced U.S. patent application Ser. No. 09/480,908 discloses a fast convolution method whose central idea (in one-dimensional embodiments) is to approximate a smooth kernel d by a polynomial spline kernel \u0192 (where \u0192 is a spline function \u0192(x) which is piecewise a polynomial of degree \u03b4 with L pieces f1(x)), and then to use appropriate operators that annihilate (or flatten) each polynomial of given degree (in a manner to be explained) to calculate the convolution of \u0192 and p quickly. In some embodiments, the smooth kernel d is approximated by a spline kernel \u0192 which is not a polynomial spline kernel, but which consists of L pieces defined over adjacent segments of its domain (in typical two-dimensional cases, the latter spline kernel is a radially symmetric function whose domain is some continuous or discrete set of values of the radial parameter). Though \u201cspline\u201d convolution as described in U.S. application Ser. No. 09/480,908 has features reminiscent of conventional wavelet schemes and is an O(N) algorithm (as are wavelet schemes), an advantage of \u201cspline\u201d convolution is that it can be performed (on data indicative of a pattern p consisting of N data values) with cN arithmetic operations (multiplications and additions), whereas conventional wavelet convolution on the same data would require bN arithmetic operations, where the factor \u201cb\u201d is typically (i.e., with typical error analysis) significantly larger than the factor \u201cc.\u201d In other words, the implied big-O constant for the spline convolution is significantly smaller than the typical such constant for conventional wavelet convolution.\nSpline convolution, as described in U.S. application Ser. No. 09/480,908, is a method for performing cyclic or acyclic convolution of a pattern \u201cp\u201d (i.e., data indicative of a pattern \u201cp\u201d) with a smooth diffusion kernel d, to generate data indicative of the convolution result r=Dp, where D is the circulant of d. The pattern \u201cp\u201d can be one-dimensional in the sense that it is determined by a continuous (or discrete) one-dimensional domain of data values (e.g., pixels), or it can be two-dimensional in the sense that it is determined by a continuous two-dimensional domain of data values (or a two-dimensional array of discrete data values), or p can have dimension greater than two. In typical discrete implementations, the pattern p is one-dimensional in the sense that it is determined by a discrete, ordered set of data values (e.g., pixels) pi, where i varies from 0 to N\u22121 (where N is the signal length), or it is two-dimensional in the sense that it is determined by an array of data values pij, where i varies from 0 to N\u22121 and j varies from 0 to N\u2032\u22121, or it has dimension greater than two (it is determined by a three or higher-dimensional set of data values). Typically, the kernel d is determined by an array of data values dij, where i varies from 0 to N\u22121 and j varies from 0 to N\u2032\u22121 (but the kernel d can alternatively be determined by a discrete set of data values d0 through dN\u22121).\nIn some embodiments described in U.S. application Ser. No. 09/480,908, the convolution Dp is accomplished by performing the steps of:\n(a) approximating the kernel d by a polynomial spline kernel \u0192 (unless the kernel d is itself a polynomial spline kernel, in which case d=\u0192 and step (a) is omitted);\n(b) calculating q=Bp=\u0394\u03b4+1Fp, where F is the circulant of kernel \u0192, and \u0394\u03b4+1 is an annihilation operator (whose form generally depends on the degree \u03b4 of the polynomial segments of \u0192) which operates on circulant F in such a manner that \u0394\u03b4+1F=B is sparse; and\n(c) backsolving \u0394\u03b4+1r=q to determiner=Fp.\nIn cases in which the kernel d is itself a polynomial spline kernel (so that d=\u0192 and F=D), the method yields an exact result (r=Dp). Otherwise, the error inherent in the method is (\u0192\u2212d)\u00d7p, where \u00d7 denotes convolution, and thus the error is bounded easily.\nIn one-dimensional cases (in which the pattern to be convolved is a one-dimensional pattern of length N), \u0394\u03b4+1 has the form of the N\u00d7N circulant matrix defined as follows:       \u0394          \u03b4      +      1        =      [                                        +                          (                                                                                          \u03b4                      +                      1                                                                                                            0                                                              )                                                            -                          (                                                                                          \u03b4                      +                      1                                                                                                            1                                                              )                                                            +                          (                                                                                          \u03b4                      +                      1                                                                                                            2                                                              )                                                            -                          (                                                                                          \u03b4                      +                      1                                                                                                            3                                                              )                                                \u22ef                          0                                      0                                      +                          (                                                                                          \u03b4                      +                      1                                                                                                            0                                                              )                                                            -                          (                                                                                          \u03b4                      +                      1                                                                                                            1                                                              )                                                            +                          (                                                                                          \u03b4                      +                      1                                                                                                            2                                                              )                                                \u22ef                          0                                      \u22ee                                      \u2003                                                \u2003                                                \u2003                                                \u2003                                    \u22ee                                                  -                          (                                                                                          \u03b4                      +                      1                                                                                                            1                                                              )                                                            +                          (                                                                                          \u03b4                      +                      1                                                                                                            2                                                              )                                                \u22ef                                      \u2003                                    0                                      +                          (                                                                                          \u03b4                      +                      1                                                                                                            0                                                              )                                            ]  in which each entry is a binomial coefficient, and \u03b4 is the maximum degree of the spline segments of spline kernel \u0192. For example, \u03b4=2 where the spline kernel \u0192 comprises quadratic segments. In two- or higher-dimensional cases, the annihilation operators can be defined as       \u0394    =                            \u2202                      x            1                                n            1                          \u2062                              \u2202                          x              2                                      n              2                                \u2062                      \u2003                    \u2062          \u2026                    \u2062              \u2003            \u2062              \u2202                  x          d                          n          d                      ,where \u2202xhnh is the n-th partial derivative with respect to the h-th of d coordinates. For example, the Laplacian       \u2207    2    \u2062      =                            \u2202                      x            1                    2                \u2062                  +                      \u2003                    \u2062          \u2026                    \u2062              \u2003            +              \u2202                  x          d                2            will annihilate piecewise-planar functions of d-dimensional arguments \u0192=\u0192(x1, x2, . . . xd).\nIn the one-dimensional case, the end points of each segment (the \u201cpivot points\u201d) of spline kernel \u0192 may be consecutive elements di and di+1 of kernel d, and step (a) can be implemented by performing curve fitting to select each segment of the spline kernel as one which adequately matches a corresponding segment of the kernel d. In some implementations, appropriate boundary conditions are satisfied at each pivot point, such as by derivative-matching or satisfying some other smoothness criterion at the pivot points.\nIn some implementations described in application Ser. No. 09/480,908, step (c) includes a preliminary \u201cignition\u201d step in which a small number of the lowest components of r=Fp are computed by exact multiplication of p by a few rows of F, and then a step of determining the rest of the components of r using a natural recurrence relation determined by the spline kernel and the operator \u0394\u03b4+1. For example, in the one-dimensional case, the lowest components of r are r0, r1, . . . , r\u03b4, where \u201c\u03b4\u201d is the maximum degree of the spline segments of spline kernel \u0192 (for example r0, r1, and r2 where the spline kernel comprises quadratic segments), and these (\u03b4+1) components are determined by exact multiplication of p by (\u03b4+1) rows of F. The (\u03b4+1) components can alternatively be determined in other ways. Then, the rest of the components \u201cr\u03b4\u201d are determined using a natural recurrence relation determined by the operator \u0394\u03b4+1. The \u201cignition\u201d operation which generates the components r0, r1, . . . , r\u03b4, can be accomplished with O(N) computations. The recurrence relation calculation can also be accomplished with O(N) computations.\nIn other embodiments, the method disclosed in U.S. application Ser. No. 09/480,908 for performing the convolution r=Dp includes the steps of:\n(a) approximating the kernel d by a polynomial spline kernel \u0192 (unless the kernel d is itself a polynomial spline kernel, in which case d=\u0192 and step (a) is omitted);\n(b) calculating q=Bp=\u0394\u03b4Fp, where F is the circulant of kernel \u0192 and \u0394\u03b4 is a flattening operator (whose form generally depends on the degree \u03b4 of the polynomial segments of F, and which operates on circulant F such that B=\u0394\u03b4F is almost everywhere a locally constant matrix); and\n(c) backsolving \u0394\u03b4r=q to determine r=Fp. In one-dimensional cases (in which p has length N), \u0394\u03b4 has the form of the N\u00d7N circulant matrix:       \u0394    \u03b4    =      [                                        +                          (                                                                    \u03b4                                                                                        0                                                              )                                                            -                          (                                                                    \u03b4                                                                                        1                                                              )                                                            +                          (                                                                    \u03b4                                                                                        2                                                              )                                                            -                          (                                                                    \u03b4                                                                                        3                                                              )                                                \u22ef                          0                                      0                                      +                          (                                                                    \u03b4                                                                                        0                                                              )                                                            -                          (                                                                    \u03b4                                                                                        1                                                              )                                                            +                          (                                                                    \u03b4                                                                                        2                                                              )                                                \u22ef                          0                                      \u22ee                                      \u2003                                                \u2003                                                \u2003                                                \u2003                                    \u22ee                                                  -                          (                                                                    \u03b4                                                                                        1                                                              )                                                            +                          (                                                                    \u03b4                                                                                        2                                                              )                                                \u22ef                                      \u2003                                    0                                      +                          (                                                                    \u03b4                                                                                        0                                                              )                                            ]  in which each entry is a binomial coefficient, and \u03b4 is the maximum degree of the spline segments of spline kernel \u0192. In higher-dimensional cases, the flattening operator \u039467  is defined similarly.\nIn other embodiments disclosed in U.S. application Ser. No. 09/480,908, the convolution Dp (where D is the circulant of smooth kernel d) includes the steps of:\n(a) approximating the kernel d by a spline kernel \u0192 which is not a polynomial spline kernel (unless the kernel d is itself such a spline kernel, other than a polynomial spline kernel, in which case d=\u0192 and step (a) is omitted);\n(b) calculating q=Bp=AFp, where F is the circulant of kernel \u0192 and A is an annihilation or flattening operator, where A operates on circulant F in such a manner that AF=B is sparse when A is an annihilation operator, and A operates on circulant F in such a manner that AF=B is almost everywhere a locally constant matrix when A is a flattening operator; and\n(c) back-solving Ar=q to determine r=Fp.\nTo better appreciate the advantages of the present invention over conventional convolution methods, we next explain two types of conventional convolution methods: Fourier-based convolution and wavelet-based convolution.\nAs is well known, Fourier-based convolution relies on the elegant fact that if F is a Fourier matrix, sayFjk=e\u22122\u03c0ijk/Nthen the transformation FDF\u22121 of the circulant is diagonal, whence we compute:Dp=F\u22121(FDF\u22121)Fp,where the far-right operation Fp is the usual Fourier transform, the operation by the parenthetical part is (by virtue of diagonality) dyadic multiplication, and the final operation F\u22121 is the inverse Fourier transform. For arbitrary D one requires actually three Fourier transforms, because the creation of the diagonal matrix FDF\u22121 requires one transform. However, if D is fixed, and transformed on a one-time basis, then subsequent convolutions Dp only require two transforms each, as is well known. The complexity then of Fourier-based cyclic convolution is thus O(N log N) operations (i.e., on the order of N log N multiplications and additions) for convolving a pattern p of length N (a pattern determined by N data values), because of the 2 or 3 FFTs (Fast Fourier Transforms) required. It should be noted that the Fourier method is an exact method (up to round-off errors depending on the FFT precision).\nAnother class of conventional convolution methods consists of wavelet convolution methods, which, by their nature, are generally inexact. The idea underlying such methods is elegant and runs as follows in the matrix-algebraic paradigm. Assume that, given an N-by-N circulant D, it is possible to find a matrix W (this is typically a compact wavelet transform) which has the properties:    (1) W is unitary (i.e. W\u22121 is the adjoint of W);    (2) W is sparse; and    (3) WDW\u22121 is sparse,where \u201csparse\u201d in the present context denotes simply that any matrix-vector product Wx, for arbitrary x, involves reduced complexity O(N), rather than say O(N2).\nWith the assumed properties, we can calculate:Dp=W\u22121(WDW\u22121)Wpby way of three sparse-matrix-vector multiplications, noting that unitarity implies the sparseness of W\u22121. Therefore the wavelet-based convolution complexity is O(N) for convolving a pattern p determined by N data values, except that it is generally impossible to find, for given circulant D, a matrix W that gives both sparsity properties rigorously. Typically, if the convolution kernel d is sufficiently smooth, then a wavelet operator W (which is sparse) an be found such that within some acceptable approximation error the property (3) above holds. Above-noted properties (1) and (2) are common at least for the family of compact wavelets (it is property (3) that is usually approximate).\nAn advantage of \u201cspline\u201d convolution (in accordance with the teaching of U.S. application Ser. No. 09/480,908) over conventional wavelet convolution is that it can be performed (on data indicative of a pattern p comprising N data values) with cN arithmetic operations, whereas conventional wavelet convolution on the same data would require bN arithmetic operations, where (assuming typical error budgets) the factor \u201cb\u201d is significantly larger than the factor \u201cc.\u201d Among the other important advantages of the \u201cspline\u201d convolution method of application Ser. No. 09/480,908 (over conventional convolution methods) are the following: spline convolution is exact with respect to the spline kernel f, whereas wavelet convolution schemes are approximate by design (and error analysis for wavelet convolution is difficult to implement); the signal lengths for signals to be convolved by spline convolution are unrestricted (i.e., they need not be powers of two as in some conventional methods, and indeed they need not have any special form); and spline convolution allows acyclic convolution without padding with zeroes.\nSeparated-spline convolution in accordance with the present invention (like spline convolution in accordance with U.S. application Ser. No. 09/480,908) is an O(N) method for convolving a pattern p determined by N data values. Separated-spline convolution in accordance with the present invention has an advantage over spline convolution in accordance with U.S. application Ser. No. 09/480,908 in that separated-spline convolution in accordance with the invention can be performed (on data indicative of a two- or higher-dimensional patterns consisting of N data values) with dN arithmetic operations (multiplications and additions), whereas spline convolution in accordance with U.S. application Ser. No. 09/480,908 on the same data would require cN arithmetic operations, where the factor \u201cc\u201d is larger, and typically significantly larger, than the factor \u201cd.\u201d In other words, the implied big-O constant for separated-spline convolution according to the present invention is significantly smaller than the typical implied big-O constant for spline convolution as described in U.S. patent application Ser. No. 09/480,908."}
-{"text": "The direct microscopy and the culture methods each have pitfalls In the past 20-25 years the direct visualization of bacteria in urine has largely been abandoned in favor of the methods involving culturing and counting the colonies of bacteria. Indeed virtually all of the studies of the significance of bacteriuria are based upon culturing the urine, and the direct microscopic examination of urine has been relegated to the status of a quick but inadequate screening procedure which may be helpful because it can be correlated with the culture methods.\nAny culture method requires that the bacteria will grow in the laboratory in the medium selected and in the time allotted. If the bacteria are damaged or dead when they left the body, then they will not grow There are many reasons why bacteria in urine would be damaged. The ionic strength or osmolality of the solution may be damaging. The oxidation potential of urine is usually too high (e.g., +0.22 to +0.25 volts). There may be a noxious metabolite in urine. (Human antibodies have been identified in urine and they have been demonstrated to be deposited on bacteria in urine.) Any or many of these factors may render a given bacterium non-viable in vitro. Finally, if the medium used is inappropriate for the growth of the particular organisms present, they will not grow. It can be readily shown by staining and microscopy that many of the bacterial forms found under the microscope were not alive at the time the specimen was obtained. For example, some do not contain nucleic acid, a biochemical component essential to life. Should all of the bacteria in a given specimen be devoid of nucleic acid, then none will grow and the culture of urine remains sterile. Indeed many urine specimens from sick patients containing huge numbers of bacteria will not yield a thriving bacterial culture in the hospital bacteriology laboratory. When the laboratory reports \"no growth\" the clinician may abandon the possibility of significant bacteriuria, and hence the possibility of an infectious cause. Nonetheless, these dead, damaged or fastidious bacteria, though they do not grow in culture, may in vivo have caused or exacerbated the patient's illness. FIG. 2 shows innumerable small cocci found, using the novel method taught herein, in the urine of a patient for whom the attempt to demonstrate bacteriuria by culture technique was completely unsuccessful. With regard to the direct examination of the urine, it must be noted that although bacteria may be seen in urine at only 100 diameters magnification, the size of the image is not the only consideration. Should the optical density and refractive index of a dead bacterium be near that of the medium, then it would not be detected by ordinary light microscopy. It may be seen by staining or perhaps by some specialized lighting. (Even then, as pointed out by Kunin, round bacteria cannot be distinguished from other near round particles such as crystals.) In my method the urine is examined wet at 100 to 400 diameters magnification, but it is also dried and prepared in a particular way so as to retain and preserve the bacterial structure through staining. In particular, I have found that urine contains lipids which act as detergents. Should they be allowed to remain on the slide when an aqueous dye is applied to the slide, then much of the sediment (including bacteria) will wash off of the slide and the preparation will be lost. This is a major reason why past attempts to study bacteriuria have failed. The photograph of FIG. 1 illustrates the difference. In that photograph one slide had been prepared in the standard way and the other has been washed with a lipid solvent. After staining much more sediment is found on the washed slide. Most of the sediment had washed off in the standard preparation. A chromatogram of the lipids removed reveals several lipids in the range of polarity of the phospholipids (e.g., lecithin, phosphatidylserine, etc. which substances act as detergents) but they do not contain appreciable phosphorus and thus they are not phospholipids. Standard methods of preparing and staining urine specimens, such as that of Melnick, U.S. Pat. No. 4,225,669, do not provide for the precautionary removal of these lipids.\nRheumatoid arthritis (RA) is a chronic inflammation of the joints, generally regarded as a systemic autoimmune disorder. Its etiology is unknown, but it has been postulated that it is associated with microbial infection. See, e.g., D. C. Demonde, ed., Infection and Immunity in the Rheumatic Diseases, 95-287 (Blackwell Scientific Publications, London: 1976). The evidence, however, until the present discovery, was inconclusive. See, e.g., D. J. McCarty, et al., ed., Arthritis and Allied Conditions: A Textbook of Rheumatology, ch 28 at 417 (9th ed. 1979); R. G. Petersdorf, et al., ed., Harrison's Principles of Internal Medicine, Part Six, Chapter 346, at 1977 (McGraw Hill: 1983). Bacteriuria has not been associated with RA, and indeed one authority remarks \"Urinary abnormalities are relatively uncommon in RA . . . Urinary tract infection was not found to be increased in RA patients.\" McCarty, supra, chapter 33, page 499, citing Ann. Rheum. Dis., 27: 345 (1968). Hypertension is a chronic elevation of blood pressure resulting from the obstruction of blood flow within the kidney (secondary hypertension) or without apparent cause (essential hypertension). One kidney disorder associated with secondary hypertension is pyelonephritis, the inflammation of the renal pelvis of the kidney as a result of bacterial infection, usually responsive to antibiotics. It has not been reported, however, that there is any correlation between essential hypertension and asymptomatic bacteriuria (bacteriuria observed in patients not reporting symptoms of urinary tract disorders). According to N. M. Kaplan, Clinical Hypertension, 14 (3d. ed. 1982), bacteriuria is found in 2-5% of hypertensives. Most of these positive cultures were of gram-negative rods. The method of the present invention has demonstrated a much higher incidence of bacteriuria in hypertensives, perhaps as high as 90%, and that cocci or \"exploded cocci\" are found in considerable numbers."}
-{"text": "1. Field of the Invention\nThe present invention generally relates to a flash memory and a method of forming a flash memory, and more specifically a method of forming a flash memory cell using an asymmetric control gate with a sidewall floating gate.\n2. Description of the Related Art\nA flash memory is unique in providing fast compact storage which is both nonvolatile and rewritable.\nIn a flash memory, the threshold voltage Vt for conduction of a field effect transistor (FET) changes state depending upon the amount of charge stored in a floating gate (FG) part of the FET. The floating gate is a charge storing region which is isolated from a more traditional gate conductor CG (control gate or \"wordline\") by a thin dielectric. The states of the Vt change with the amount of charge stored by the FG.\nSince the FG directly controls conductivity between source and drain in a channel, the state of a FG memory cell is determined by applying certain voltages to the source or drain of the FET and observing whether the FET conducts any current.\nFlash memory cells with a sidewall floating gate occupy a smaller area than those with conventional (layered) floating gates. For example, in U.S. Pat. No. 5,115,288, sidewall gates were formed on one side of the wordline by employing an extra mask. Sidewall spacers were formed on both edges of the wordline, then removed along one of the edges using the extra mask and an etching operation. Thus, the conventional approach uses a trim mask to define the floating gate.\nHowever, this approach is expensive and requires good control of the overlay for the spacer removal mask.\nOther conventional structures also are known. For example, in one structure, polysilicon spacers formed on both sidewalls are used for the floating gate. One spacer sits on top of the tunnel oxide area for programming. The other spacer is called an \"added-on floating gate\". Both spacers are linked by a polysilicon body. However, a large cell size results.\nIn a second conventional structure, only one polysilicon spacer is used as the floating gate. A mask must be aligned to the top of the control gate, to remove the other floating gate spacer. Hence, the control gate cannot be small, since, otherwise, any misalignment will cause a problem. Therefore, this cell has difficulty in being down-scaled.\nIn yet another conventional structure, similar to the second conventional structure described above, a mask is needed to remove a sidewall spacer floating gate. Further, this spacer has a re-entrant corner which is very difficult to be completely removed.\nThus, the conventional methods require extra process steps, material and more precise lithographic alignment, thereby resulting in increased manufacturing costs."}
-{"text": "The present invention relates generally to a circuit arrangement used as an interface for a sensor, in particular, for a pumped reference oxygen sensor used in combination with a combustion engine.\nOxygen sensors are particularly used in combination with combustion engines using controlled catalytic converters. The use of catalytic converters in cars started much earlier in the USA in comparison with Europe. However, innovation cycles are often slower in the US automotive technology. Therefore, often older technology is used for longer periods in the USA as compared to European countries. For example, the binary lambda oxygen sensor which is used to regulate the gasoline mixture in a combustion engine, comprises the negative terminal being electrically coupled with the sensor housing in embodiments of the first generation. The reasons for this connection relates to an easier and cheaper construction. However, this construction also results in an electrical coupling between the sensor housing and the engine through the fixture of the sensor within the muffler arrangement. This connection is disadvantageous because shifted ground potentials within a motor vehicle. For example, the engine ground is usually more negative than the ground of the motor control unit. This potential shift is due to switching of high load currents, e. g., 10 . . . 15A, and the inner resistance of the ground wires, e. g. 20 . . . 40 m\u03a9. Typical voltage shifts are in the range of \u2212300 . . . +600 mV. In addition, any on/off switching of high loads causes voltage overlay peaks of up to multiple 100 m Vss. These ground distortions and overlay voltages can cause serious problems with respect to evaluation of the respective sensor signals and may render the actual sensor signals completely useless.\nOxygen sensors according to newer technology are, therefore, fully isolated and, thus, avoid any electrical coupling with the engine. However, oxygen sensors of the first generation are still widely used, in particular, in the United States because of their lower manufacturing costs. Therefore, modem motor control units must be able to interface with these kind of sensors which are not fully isolated. To this end, specific interface circuits used to be available which allowed for the evaluation of sensors whose housing is electrically coupled with the engine and thus with the motor vehicle ground.\nFIG. 4 shows an example of such a known circuit. A sensor 400 which is connected with the engine ground is coupled through resistors 430 and 415 with the input terminals of the interface circuit proper. Capacitors 425 and 435 are coupled between the input terminals of the interface circuit and the interface ground. One of the sensor connections is furthermore coupled through resistor 410 with a supply voltage Vcc and through capacitor 420 with the interface ground. The interface circuit comprises a first switch 440 which is coupled with the input terminals. The switch output is coupled through a capacitor 445 with the input of a second switch 460 and through capacitor 450 with the input of a third switch 455. The first output of switch 455 is connected with the supply voltage Vcc. The second output of switch 455 and the first output of switch 460 are coupled with the interface ground. The second output of switch 460 is connected to the inverting input of an integrator consisting of operational amplifier 475 and capacitor 470 in its feedback loop. A fourth switch couples the inverting input of the integrator with either the supply voltage or the output of the integrator. The output of the integrator is coupled with the first input of a fifth switch 495 which is controlled by the output signal of a comparator 405. The second input of switch 495 is coupled through resistor 490 with the supply voltage. The first input of comparator 405 is coupled with the second interface input terminal and the second input of comparator 405 receives a voltage signal being equal to half the supply voltage. Furthermore, a timing circuit is provided which generates control signals for switches 440, 455, 460, and 480.\nAll switches are implemented as CMOS switches. Capacitor 445 is used as a transfer capacitor for eliminating the common mode of the input signal. To this end, the capacitor is switched in a first position between interface ground and the second input terminal and in a second position between the inverting input of operational amplifier 475 and the first interface input terminal with a high frequency. Thus, the CMOS switch operates like a resistor. Capacitor 465 operates as a feedback capacitor in a similar way. These two capacitors operating as resistors form together with the operational amplifier/integrator an inverting amplifier. The bias capacitor 450 together with CMOS switch 455 are used to generate a small bias current which is fed to the sensor 400 and which will not influence the measurement when the sensor is in operating mode, i.e. the sensor has low resistance.\nOne of the disadvantages of this circuit arrangement is that the CMOS switches at the input of the circuit must comply with a high standard. This renders this circuit expensive and interference-prone. In addition, this circuit must withstand the required negative input voltages. thus, additional protective measurements, such as, isolation and charge pumps (not shown) must be provided. Furthermore, the CMOS switches must be able to tolerate a relatively high input voltage of up to 12V in case of a short circuit of the sensor. This is particularly difficult because the supply voltage is usually only 5V. Integrated circuits using this technology need furthermore additional isolation/separation measurements if more than one interface circuit is provided to prevent any cross over influence of the channels and to prevent a latch-up.\nThe bias current generated by switch 455 and capacitor 450 is used to detect a connection failure between the first input terminal and the sensor. In such a case, the bias current will overdrive the operational amplifier. A similar scenario takes place in case of a short circuit between the first input terminal and the positive terminal of the battery. The output voltage in both cases will be approximately 0V. To detect any interruption between the second input terminal and the sensor additional circuitry is necessary. This additional circuitry is shown in FIG. 4 with resistor 410, capacitor 420, comparator 405 and CMOS switch 495. During normal operation the current generated by resistor 410 will flow to the engine ground through the electrical coupling of the housing of sensor 400 and will not influence the measurement. However, in case of an interruption of this connection the potential at the second input terminal will raise to the supply voltage, for example, 5V. In case of a normal operational temperature of the sensor, i.e. low resistance of the sensor, the operational amplifier will be driven to its positive limit, e.g., 5V. However, in case of a cold sensor (during the start up phase of the engine) the sensor will have a high resistance and the internally generated bias current will put the circuit into an undefined state. To prevent such a state, comparator 405 will compare the potential at the second input terminal with Vcc/2. If the potential is above this threshold, comparator 405 will control switch 490 to select a constant output voltage to signalize this error.\nAs described above, the prior art interface circuit is highly cumbersome and requires additional evaluation of the generated output signal. Furthermore, this type of interface circuit is not in production anymore and, thus, not available for new construction which specifies the use of a non-isolated sensor."}
-{"text": "The invention relates to manufacturing an inner conductor of a resonator.\nResonator structures of a high frequency area, a radio frequency area in particular, are used e.g. in base stations of mobile telephone networks. Filters may utilize resonator structures e.g. as adapting and filtering circuits in transmitter and receiver units of the base stations.\nA resonator structure comprises an inner conductor of the resonator attached to an attachment surface, which in practice most often is an end, such as a bottom or a cover, of a housing structure serving as an outer conductor of the resonator structure. The inner conductor is thus short-circuited to the attachment surface, i.e. in practice to the outer conductor. A short-circuited end of the inner conductor, at which the inner conductor is thus short-circuited to the outer conductor, is also called an inductive end owing to the fact that signal coupling at the short-circuited end is mainly carried out inductively.\nAt a second end of the inner conductor, the inner conductor is galvanically separated from the outer conductor, so this end is the xe2x80x9cfreexe2x80x9d end of the inner conductor. The free end of the inner conductor is also called a capacitive end of the inner conductor owing to the fact that signal coupling at this end is mainly carried out capacitively. The outer conductor and the inner conductor located within a section defined by the outer conductor together form a resonance circuit. In practice, the resonator structures often comprise a plurality of circuits, i.e. the resonator structure comprises several pairs comprising an inner conductor and an outer conductor, i.e. each section formed by the outer conductor comprises a separate inner conductor. The resonance circuits of a multi-circuit resonator structure together form a desired frequency response for the resonator structure.\nNormally in a coaxial resonator, the inner conductor of the resonator is a straight wire or a pin attached only to the bottom of the resonator. Such a resonator is long and thus takes a lot of space. The resonator pin is quite easy to manufacture. The problem then is, however, how to adjust the coupling of the resonator since it is difficult to attach such a controlling element to the resonator pin that would enable the resonator to be easily coupled to e.g. an adjacent resonator. Furthermore, the capacitive coupling provided by the wire-like inner conductor is poor.\nIn order to decrease the space required by the resonator, for instance a helix coil is used as the inner conductor, in which helix coil the same operational length fits into a shorter space since the resonator in the helix resonator is formed as a coil. The helix coil is, however, difficult to manufacture. A further drawback is that it is extremely difficult to attach to the helix coil a coupling wire or other such projection necessary when the coupling between two resonance circuits is to be adjusted. A further problem with the helix resonators is the difficulty to support them and carry out the temperature compensation. An inner conductor implemented by utilizing a helix coil cannot provide a high-quality capacitive coupling.\nA known solution for controlling the resonance frequency of a resonator circuit is a solution wherein an adjuster bolt located in the cover of a filter serves as the frequency controlling element, and the distance of the adjuster bolt with respect to the free end of the resonator located in a section under the cover is adjusted by turning the bolt. The solution is not the best possible one since it requires additional structures on the outer surface of the housing. A further problem is that the adjuster bolt requires that the cover of the filter should be thick or the cover should at least comprise a thicker section to enable threads to be provided on the cover for the adjuster bolt, or, alternatively, to enable a nut-like part with threads attached to the cover to be used. The cover has to be thick particularly because it also needs to be rigid in order to prevent the distance of the frequency controlling element in the cover with respect to the resonator from changing after the controlling procedure and from further causing the capacitance, and thus the resonance frequency, to change in an undesired manner.\nAn object of the invention is thus to provide a method of manufacturing an inner conductor of a resonator, and an inner conductor so as to enable the above-mentioned problems to be alleviated. This is achieved by a method disclosed in the introduction, characterized by manufacturing at least part of the inner conductor from a uniform, electrically conductive material blank by utilizing a deep-drawing method wherein the blank is struck or pressed with a tip of an impact device, whereby during each stroke or pressing, the tip draws more and more blank material in the direction of the stroke.\nThe invention further relates to an inner conductor of a resonator comprising a first end and a second end, which is free.\nThe inner conductor of the invention is characterized in that at least part of the inner conductor is deep-drawn from a uniform, electrically conductive blank.\nPreferred embodiments of the invention are disclosed in the dependent claims.\nThe idea underlying the invention is that the inner conductor is manufactured by utilizing a deep-drawing method.\nSeveral advantages are achieved by the method and inner conductor of the invention. The deep-drawing method enables the inner conductor and a flange located at the free end thereof to be manufactured virtually simultaneously. In addition, a potential projection or a site for the same can be manufactured in connection with manufacturing the inner conductor. The drawing method is a quick and low-cost way to manufacture inner conductors. The drawing method enables flanges and projections for the inner conductors to be manufactured that are all integrated in the same uniform material piece. Therefore, the inner conductor is mechanically strong.\nSince the inner conductor is deep-drawn, the surface of the inner conductor is extremely smooth, which enables the inner conductor to be readily coated e.g. with silver. Thanks to the smoothness of the surface, the surface area to be coated is smaller than it would be if the surface was uneven. It thus takes less coating material to coat an even surface than an uneven one.\nAn inner conductor manufactured by utilizing the deep drawing method has a small surface resistance, so the electric loss of the resonator remains small and the Q factor of the resonator can be retained good.\nA further advantage of the deep drawing method is that the inner conductor can be manufactured e.g. from a copper blank, in which case the resulting inner conductor does not necessarily have to be coated. The inner conductor manufactured of copper is attached by a specific sleeve, which means that the inner conductor made of copper does not have to be mechanized for the screw threads in a fixing screw.\nSince it is possible to attach the inner conductor by a sleeve, the thickness of the walls of the inner conductor can be retained small, which gives a lightweight inner conductor. The advantage provided by the inner conductor being light is that it is highly tolerant e.g. of vibration. Consequently, external vibration does not easily cause the inner conductor to move or become detached. The structure and attachment of the inner conductor thus enable intermodulation noise to be reduced."}
-{"text": "A typical multi-color dye donor web that is used in a dye transfer or thermal printer is substantially thin and has a repeating series of three different rectangular-shaped color sections or patches such as a yellow color section, a magenta color section and a cyan color section. In addition, there may be a transparent colorless laminating section immediately after the cyan color section.\nEach color section of the dye donor web consists of a dye transfer area which is used for dye transfer printing and a pair of opposite longitudinal edge areas alongside the dye transfer area which often are not used for printing. The dye transfer area may be about 152 mm wide and the two longitudinal edge areas may each be about 5.5 mm wide, so that the total web width is approximately 163 mm.\nTo make a multi-color image print using a thermal printer, a motorized donor web take-up spool draws a longitudinal portion of the dye donor web off a donor web supply spool in order to successively move an unused single series of yellow, magenta and cyan color sections over a stationary liner array (bead) of selectively heated resistive elements on a thermal print head between the supply and take-up spools. Respective color dyes within the yellow, magenta and cyan color sections are successively heat-transferred line-by-line, via the selectively heated resistive elements, onto a dye receiver medium such as a paper or transparency sheet or roll, to form the color image print. The selectively heated resistive elements often extend across the entire width of a color section, i.e. across the dye transfer area and the two longitudinal edge areas comprising that color section. However, only those resistive elements that contact the dye transfer area are selectively heated. Those resistive elements that contact the two longitudinal edge areas are not heated. Consequently, the dye transfer occurs from the dye transfer area to the dye receiver medium, but not from the two longitudinal edge areas to the dye receiver medium.\nAs each color section is drawn over the selectively heated resistive elements, it is subjected to a longitudinal tension particularly by the forward pulling force of the motorized donor web take-up spool. Since the dye transfer area in the color section is heated by the resistive elements, but the two longitudinal edge areas alongside the dye transfer area are not, the dye transfer area is significantly weakened and therefore vulnerable to stretching as compared to the two longitudinal edge areas. Consequently, the longitudinal tension will stretch the dye transfer area relative to the two longitudinal edge areas. This stretching causes the dye transfer area to become thinner than the non-stretched edge areas, which in turn causes some creases or wrinkles to develop in the dye transfer area, most acutely in those regions of the dye transfer area that are close to the non-stretched longitudinal edge areas. The creases or wrinkles occur most acutely in the regions of the dye transfer area that are close to the non-stretched edge areas because of the sharp, i.e. abrupt, transition between the stretched (thinner) transfer area and the non-stretched (thicker) edge areas.\nAs the dye donor web is pulled by the motorized donor web take-up spool over the selectively heated resistive elements, the creases or wrinkles tend to spread from a trailing (rear) end portion of a used dye transfer area at least to a leading (front) end portion of the next dye transfer area to be used. A known problem that can result is that the creases in the leading (front) end portion of the next dye transfer area to be used will cause undesirable line artifacts to be printed on a leading (front) end portion of the dye receiver medium. The line artifacts printed on the dye receiver medium, although they may be relatively short, are quite visible.\nThe question presented therefore is how to solve the problem of the creases or wrinkles being created in an unused transfer area so that no line artifacts are printed on the dye receiver medium during the dye transfer.\nThe Cross-Referenced Applications and Patent\nThe cross-referenced applications each disclose a thermal printer capable of preventing crease formation in a dye transfer area of a dye donor web that can cause line artifacts to be printed on a dye receiver during the dye transfer from the dye transfer area to the dye receiver. To prevent crease formation, there is provided a crease-preventing platen roller that is movable to hold a dye transfer area and the two longitudinal edge areas alongside the dye transfer area against a print head. The crease-preventing roller has a pair of roller end portions that apply a constant pressure against the two longitudinal edge areas, and a roller main portion between the roller end portions that applies a lesser pressure against the dye transfer area. Since the pressure applied against the two longitudinal edge areas is greater than the pressure applied against the dye transfer area, the two longitudinal edge areas will be stretched the same as the dye transfer area, so that creases will not be formed in the dye transfer area. This is so notwithstanding that the dye transfer area is heated by the print head, but the two longitudinal edge areas are not.\nIn contrast to the cross-referenced applications, the referenced incorporated (prior art) patent discloses a thermal printer that is adapted to optimize print image quality by preventing undesired pressure variations along the line of contact between the dye donor web and the linear array of selectively heated resistive elements in the thermal print head (the patent does not discuss the problem of crease formation). To optimize print image quality, there is provided a plurality of pressure applying rods that bear down on the thermal print head to urge the selected heated resistive elements into pressure contact with the dye donor web. The amount of pressure applied by each rod and the location of each rod along the print head is individually adjusted in response to sensed changes in different operating parameters that negatively affect print image quality, such as print head temperature when printing dark vs. light image portions, and thickness and/or stiffness of the dye donor web."}
-{"text": "1. Field of the Invention\nThe present invention relates to a hologram memory apparatus for recording and reconstructing spatial pattern information. More precisely, the present invention relates to an improvement of such an apparatus of which the structure is greatly simplified by using the one-dimensional hologram recording.\n2. Description of the Prior Art\nAs a conventional method and apparatus for recording information in a one-dimensional hologram, it is known to use a computer-generated hologram, in which an input information is Fourier transformed by means of a computer to produce an electric holographic signal and an electron beam or a laser beam is modulated by the electric signal and is used for recording. For instance, Adam Kozma et al had submitted a report in the 1971 IEEE/OSA Conference on Laser Engineering and Applications (Session No. 152). This suggested system has an advantage in that a spatial modulator is not required; however, it contains a disadvantage in that the writing speed cannot be made high enough owing to a long computation time.\nThere are various disclosures with respect to the two-dimensional page-oriented hologram memory. This kind of apparatus, however, had encountered a number of technical difficulties mainly due to the complicated construction of the structural devices since it requires various elements such as a two-dimensional spatial modulator, i.e. page composer or a two-dimensional photo detector matrix, etc.\nFurthermore, there had been an attempt to realize a simplified hologram memory system considering applications as wideband recorders. For instance, Andrew Bardos disclosed such a system in Applied Optics, Vol. 13, No. 4, April 1974, pp. 832-840. In this system, a time series of a given electrical pulse information is rearranged into multi-channel information and by using such multi-channel information the respective channel of a multi-channel acousto-optic light deflector is controlled as a spatial modulator and thus the input signal light beam is formed. The input signal light beam is subjected with a doppler shift. Accordingly, in the recording of the hologram in order to stabilize the interference fringe, the reference light should also be subjected with such a doppler shift to the same extent with the signal light by using another acousto-optic light deflector, names an acoustic beam splitter, to give diffractive deflection for the laser beam. Accordingly, the optical path for the reference light beam is separated from that of the signal light path at the acoustic beam splitter and thus the strip-shaped hologram recording is made through a writing lens after readjustment of the signal light beam and the reference light beam to be parallel with each other. There exist two problems in this step. One is the stability of the system due to use of a separate optical path. Another is the demand for large aperture lens systems, which make them costly and difficult to produce."}
-{"text": "1. Field of the Invention\nThis invention relates to improved apparatus for error correction of digital data, and in particular to correction of digital data representing a pictorial image.\n2. Description Relative to the Prior Art\nWhile much digital data, such as that derived from computation and general data processing has no redundancy, data generated by digital processing of images is characterized by substantial amounts of redundant information. The established method of converting images into digital data entails projecting the image onto a sensor having an analog output, and then scanning the sensor and digitizing its output. The redundancy inherent in the image is carried over into the resultant digital data stream; the digital values themselves exhibiting the redundancy of the original pictorial image.\nReferring to FIG. 1, a scene 10 is scanned in horizontal segments by an optical scanning device 12 whose output is a raster 14 containing the pictorial information in the form of electrical analog signals. The scanning device 12 may be any one of the optical scanners known in the art; for example, the laser scanner, kinescope scanner, or CCD sensor and associated scanning circuits. The electrical signals are then sequentially digitized by means of a sampling analog to digital converter 16 whose output is the digital representation of the image intensity on a line by line basis. Such image samples are conventionally referred to as \"pixels\". Generally, the digitizer 16 output comprises the pixel value expressed in binary form wherein eight bits per sample correspond to a 256 level gray scale of image intensity. A digital eight bit sample is usually designated as a \"byte\", however, a byte may be defined to contain any discrete number of bits. This binary coded representation of the digital values of the samples is called \"pulse code modulation\", or PCM. Usually the image is digitized to allow either electromagnetic transmission or magnetic recording of the PCM encoded samples.\nIn FIG. 1, the line 18 represents a left-to-right scan across the \"sky\" portion 19 of the scene 10, and it will be appreciated that the image intensity across such a scan will be essentially constant. Similarly, the intensity of the scan of the line 20 across the sky remains essentially constant until the edge of the \"cloud\" 22 is encountered, and then the image intensity changes to another essentially constant value for the duration of the scan across the \"cloud\". Upon completion of the scan of the cloud, the intensity and the sample values then revert to the values representative of the \"sky\" 19. Such scans result in sample values that are highly redundant because there is little detail in the portions of the scene being sampled. On the other hand, the detailed portions of the scene which consist of rapidly changing intensities result in samples of limited redundancy; during the scan 23 across the \"tree\" 25, rapidly varying sample values of intensity are generated. Thus, the data is characterized by non-redundant bytes derived from scans of detailed portions of the image such as the \"tree\" 25, and by redundant byte values derived from scans of the \"sky\" 19 or the \"cloud\" 22 where the intensity changes little, if at all.\nThe present inventor's copending application, Ser. No. 023,327, U.S. Pat. No. 4,761,782 focuses attention on the use of image redundancy for error correction after the transmission or storage of the digitized image data. Aberrations such as fading in an electro-magnetic transmission channel, or dropouts in a magnetic tape storage channel cause signal loss and upon recovery of the data, attendant errors. These errors result in easily discernible, detrimental streaking in the recovered image. The teaching of the above copending application ameliorates this problem by utilizing the available image redundancy as the basis for error correction in the recovered image data. It teaches the grouping of the digitized bytes representative of an image into data blocks whereby the redundant information present in the image becomes replicated in the data block and results in some of the bytes having essentially the same, or nearly the same, digital values as other bytes within the same block. Prior to transmission or storage, each block is analyzed for redundant sample values, i.e. values which are either identical or which only vary within predetermined boundaries, and check bits are accordingly affixed to the block. The configuration of the check bits defines the nature and the extent of the redundancy in the bytes comprising the block. The block containing data and check bits is then transmitted or stored. The recovered data, which has been subject to induced errors during transmission or storage, is analyzed to ascertain whether the bytes comprising the block still contain the same redundancy specified by the affixed check bits. If an error has occurred, the resultant redundancy of the received data will generally be reduced; at least one of the recovered pixels in the block will be different from what it originally was due to the error. The above mentioned copending application teaches re-encoding the received data to derive new check bits based on the received data which are then compared to the original check bits as received. This comparison generates a \"syndrome\" having the value \"0\" if the check bits derived from the received data are the same as the original check bits, and having the value \"1\" if the re-encoded check bits do not agree with the original check bits. A syndrome of value 1 flags an error and the fact that the data in the block has changed during transmission. As described in the copending application, if an error has occurred, the received data bytes are then compared among themselves to ascertain which data pattern is characteristic of the majority of received bytes. This majority pattern is assumed to be the correct pattern for all the bytes of the block including the bytes in error. The bytes in error are then corrected in accordance with the instructions inherent in the configuration of the original check bits. The correction restores the data patterns to their original configuration, and, attendantly, restores the amount of redundancy that was originally present in the block.\nThe check bits are not derived from the actual values of the data bits as is characteristic of many codes known in the art. The check bits reflect the fact that \"patterns\" of data bits within the block are the same, but the check bits do not explicitly define the structure of these patterns.\nIn considering the teaching of the prior art, it is advantageous to provide a listing of various check bit symbols used in describing the technique.\nSx=check bit x derived from data before transmission, PA1 SxR=received value of check bit Sx after transmission PA1 Sx*=re-encoded value of Sx derived from data as received after transmission PA1 SCx=value of syndrome calculated from SxR and Sx*, where for m check bits PA1 x=0,1,2 . . . m-1. PA1 127=01111111 PA1 128=10000000 PA1 132=10000100 PA1 129=10000001 PA1 127=01111111.\nA prior art embodiment employs a two bit check pattern, S1, S0, which is attached to each block of n pixels to provide an 8n+2 bit block. Check bits are assigned in accordance with the following rules:\nTABLE I ______________________________________ S1 S0 CONDITION ______________________________________ 1 1 Three most significant bits of all n pixels form indentical patterns 1 0 Only the two most significant bits of all n pixels form identical patterns 0 1 Only the most significant bits of all n pixels are the same 0 0 None of the above ______________________________________\nThe above assignment of check bits reflects the degree of redundancy carried by the three most significant bits of the n pixels comprising a block. For example, the check bits 1,1 indicate that the n pixels have the highest redundancy in their three most significant bits, i.e. the three most significant bit patterns for all n pixels are the same. On the other hand, the check bits 0,0 indicate that there is no redundancy among the most significant bits of the n pixels, i.e. all of the patterns are different. The intermediate cases of 1,0 and 0,1 reflect in-between conditions of redundancy.\nIt will be appreciated that the prior art technique depends upon the image redundancy as expressed in the three most significant binary bits of the pixel binary representation. Because of the structure of binary representation, it is possible that certain values of redundancy will fail to encode in a manner correctly representing the redundancy value. For example, a block may consist of digitized image values of 127, 128, 132, 129, 127. Encoded into an eight bit binary representation, these values are:\nWhile it is clear by consideration of the actual magnitudes of these pixels that the information is highly redundant, neither the three most significant bits of the binary values of the five pixels are the same, nor are the two most significant bits, nor are the most significant bits themselves, the same. Due solely to the structure of binary representation, these five pixels appear to be completely non-redundant rather than highly redundant. This occurs at the transition between the numbers 127 and 128 in binary representation because all the bits switch in value with the result that values just above the transition point and values just below it have no similarity in the high order bit positions. A similar \"boundary\" problem also occurs for pixels having digitized binary values in the region of 256, and 64. It will be appreciated that because of this characteristic large undetected errors may occur at these boundary points. For example, it is clear that if the most significant bit of the binary number 128 is subject to error and changes from 1 to 0, the resultant pixel value instead of being 128 is 0; a level change of 50 percent of full scale, (i.e. 256), that is not corrected in the prior art due to coding failure at the boundary.\nAn analysis of typical images containing redundancy shows that from 7 to 15 percent of blocks may fail to incorporate a block's true redundancy value when utilizing the techniques of the prior art. This is partially due to the redundancy characteristics of the image and partially due to the above described boundary problem. The mere occurrence of pixel values which straddle a boundary region is not sufficient in itself to cause errors in the reproduced image. It is necessary that a transmission error or, say, a magnetic tape storage dropout error simultaneously occur with the incidence of the pixel value in the boundary region. Because the joint probability of these two conditions occurring simultaneously is very small, acceptable error correction in many applications is provided by the technique of the prior art. For the processing of the highest quality images wherein the occurrence of any large pixel errors, even if rare, is unacceptable, additional correcting power is required. The present invention provides that power by completely eliminating the source of the errors due to the boundary problem."}
-{"text": "The present invention is directed to a gas laser of the type having a vacuum-tight sealed housing which has a capillary extending at least partially in the housing and containing a discharge channel. Magnets are located in the proximity of the capillary, the magnetic field of the magnets penetrating the discharge channel in the capillary at least over a portion of its length. Such a prior art gas laser is disclosed in U.S. Pat. No. 4,035,741. In the prior art gas laser, magnets are located outside a housing which is sealed so that the housing is vacuum-tight and are located in the proximity of a capillary that projects over the housing. A relatively complex adjustment mechanism is required in order to adjust the magnets relative to the capillary. The arrangement disclosed basically covers only one part of the capillary.\nIBM Technical Disclosure Bulletin, Vol. 22, No. 8A, January, 1980 discloses a similar arrangement in which a row of magnets is located in proximity to the discharge tube."}
-{"text": "The present invention relates to an improved method and apparatus for interconnecting large scale integrated (LSI) circuits and, more particularly, associative techniques for interconnecting LSI elements on a wafer to enable wafer scale integration.\nPresent methods of semiconductor fabrication typically require that a plurality of identical semiconductor elements or circuits be fabricated on a single wafer substrate by a series of process steps. At the conclusion of this fabrication process, the elements or circuits are tested and the defective elements identified. The wafer is then scored and diced into individual parts, each containing a complete circuit. Finally, the operable parts are packaged to provide external electrical connections and appropriate environmental protection.\nIt has long been realized that circuit costs as well as space and power requirements could be reduced if the operable parts or elements on the wafer could be interconnected on the wafer itself. This would reduce the necessity for circuit boards and interconnection wiring between elements, which would, in turn, result in a decrease in space and cost.\nA variety of methods for interconnecting elements on a wafer have been suggested. One method is to add another metalization step after the circuits or elements of the wafer are tested to interconnect the operable element on the wafer. The difficulty with this method is that a different mask is required for each individual wafer, greatly increasing the cost of the finished product.\nAnother method is the inverse of the discretionary wiring concept attempted in the 1960's to achieve LSI circuits with small scale integrated and medium scale integrated circuits. Usually, the arrays to be connected are storage arrays. Each array consists of a storage portion, an address mechanism for accessing the data and a permanent disconnection mechanism. The permanent disconnection mechanism is accomplished by blowing a fuse, or charging up a floating gate MOS (FAMOS) device, or laser burnout, etc. The disconnection isolates a defective array from the wafer bus so that the defective arrays do not interact adversely with the operative arrays. Although this method is suitable for wafer scale integration, it is limited in application. It requires a hardwired decoder and does not permit reconnection of the spring of arrays for increased reliability.\nAnother wafer integration scheme, as described in U.S. Pat. No. 3,940,740, provides for spare rows and columns on the wafer matrix for appropriate sparing in case of defective elements. One disadvantage of this configuration is that an entire row of elements must be provided, in the worst case, to spare a single defective matrix element. A preferred method would be one where spare elements can be inserted into the matrix on a random basis to allow all spare elements to be used.\nIn addition, all of these methods require extensive decoding circuits and uniquely dedicated decoding bus lines for element enable, which reduce the usable area of the wafer and add possible failure sites. Because of these reasons, wafer scale integration is seldom achieved.\nWhat is required is a wafer scale integration configuration wherein a relatively small number of spare elements and a relatively simple decoding circuit are employed to obtain high reliability and high functional circuit density at the wafer level.\nIt is also an object of this invention to allow the replacing circuit element to assume the address of the replaced element so that the system need not be reprogrammed after the replacement."}
-{"text": "A typical bicycle wheel is formed completely of metal, and includes a cylindrical hub connected by numerous metal spokes to a rim, and an axle assembly that may include a brake lying within the hub. Improved performance of bicycle wheels has been obtained by using molded plastic wheels formed of very high strength plastic material. Design of the wheel for injection molding results in a different type of wheel, wherein a limited number of heavy duty spokes are preferably utilized. U.S. Pat. No. Des. 244,272 by Raudman and Hoffman show a wheel of this type. While high performance plastics have many advantages, one important disadvantage is that they are easily damaged by heat, such as that which can be generated by a bicycle brake when repeatedly applied. Also, the limited number of heavy duty spokes in a molded plastic wheel, can result in poor absorption of shocks when applied along the axis of one of the heavy duty spokes. The limited number of spokes also can have disadvantages in providing a large open space into which a child's foot may be inserted. A molded plastic vehicle wheel, especially one designed for use on bicycles, which helps minimize the transmission of shock between the rim of the wheel and the bicycle frame, and which minimizes the possibility of heat damage to the wheel by a brake assembly mounted in the hub of the wheel, would be of considerable value in the production of bicycles as well as other vehicles."}
-{"text": "The present invention relates to a disk player which regenerates information from a disk-type recording medium and/or records information on the same, and more particularly to a disk player of a kind, which is improved in relationship between opening/closing timing of a shutter arranged on a disk cartridge for accommodating the disktype recording medium, and an information regenerating location and/or an information recording location. The present invention also relates to a disk player of a kind, which is improved in a driving mechanism thereof, for upward and downward shifting a recording mechanism which records information on the disk-type recording medium.\nAs is widely known, a disk cartridge of a disk player, which accommodates therein a disk-type recording medium, is provided with a slidable shutter for properly exposing an information surface of the disk-type recording medium. The shutter is closed normally, to thereby protect the disk-type recording medium in the disk cartridge from an external environment.\nTherefore, in conventional disk players, when information is regenerated from the information surface of the disk-type recording medium or the information is recorded on the same, the following complicated operations are required: First, the disk cartridge is carried to a predetermined location by a carrier mechanism, and at the same time the shutter of the disk cartridge is fully opened by means of a shutter opening/closing mechanism mounted in the disk player, to thereby expose the information surface of the disk-type recording medium. Then, the information is regenerated from the exposed information surface by means of an optical pickup provided in the disk player, or the information is recorded on-,the exposed information surface by means of the optical pickup and a magnetic head.\nThe above-mentioned conventional disk players are operated in the following two modes with respect to information regeneration and information recording:\nIn the first mode, when the disk cartridge is located at an information regenerating location in the disk player, the shutter is fully opened by means of the shutter opening/closing mechanism. Thereafter, while the disk cartridge is further shifted from the location to an information recording location in the disk player, the magnetic head is disposed in sliding contact with the information surface of the disk-type recording medium, to thereby carry out information recording.\nIn the first mode, however, when the disk cartridge is further shifted from the information regenerating location to the information recording location, an additional configuration for shifting the shutter opening/closing mechanism of the disk player together with the disk cartridge, or alternatively a configuration for detaching the shutter opening/closing mechanism from the disk cartridge in the above state while maintaining the fully open state of the shutter is required. Therefore, the configuration for shifting both the carrier mechanism and the shutter opening/closing mechanism or the configuration for separating the two mechanisms from each other complicates and upsizes the disk player, resulting in an increased number of component elements and hence a rise in manufacturing costs.\non the other hand, in the second mode of the disk player, the information regenerating location and the information recording location are set to the same location, and the shutter is fully opened at the location, followed by disposing the magnetic head in sliding contact with the information surface of the disk-type recording medium to carry out information recording.\nIn this mode, however, the carrier mechanism must be detached from the disk cartridge, and then the magnetic head is driven to be disposed in sliding contact with the information surface of the disk-type information recording medium, by the carrier mechanism. Alternatively, a magnetic head driving mechanism for driving the magnetic head must be additionally installed. Therefore, if the carrier mechanism is further shifted or the magnetic head driving mechanism for driving the magnetic head is additionally employed, the construction of the disk player is further complicated and upsized. In addition, the number of component elements is increased, which inevitably leads to increase manufacturing costs of the disk player.\nBesides, in either mode, to hold the information recording medium inserted into the disk player, on a turn table, a supporting member for supporting the turn table need to be relatively shifted with respect to the information recording medium, and further, during information recording, the magnetic head need to be relatively shifted with respect to the information recording medium. Thus, both the supporting member and the magnetic head need to be relatively shifted with respect to the information recording medium, which complicates and upsizes the driving mechanism, resulting in a further inconvenience, e.g. that the accuracy of a vertical location cannot be improved."}
-{"text": "The present invention relates to compounds which neutralize axillary malodor.\nVarious approaches have been taken to the problem of axillary malodor. One approach has been the use of deodorants which may contain germicides and/or a fragrance. Germicides inhibit the reproduction of the bacteria which are believed to contribute to the production of axillary malodor. Fragrances are used to mask any odor produced. Another approach to the prevention of malodor is the use of antiperspirants. Antiperspirants inhibit bacteria and reduce the amount of sweat production thereby limiting the formation of substrate which gives rise to axillary malodor.\nCertain metals and metal salts have been suggested for use in the control of odors. For example, Japanese Disclosure No. 83-222011 describes a cream for removing odor from the axilla which includes copper powder and perfume in a base of cosmetic cream. The disclosure postulates that the copper powder acts on the secretion of the odor, thereby suppressing generation of the malodor.\nAnother description of the use of metal salts against odors appears in French Patent No. 1,394,875. This French patent describes the use of water-soluble iron and copper compounds to reduce or eliminate body odors. Various salts and complexes of cupric copper and ferrous iron are mentioned, including sulfates, chlorides, acetates, gluconates, citrates, tartrates, salts of ethylenediamine tetraacetic acid sodium or potassium, and iron and copper chlorophyllins. Various iron and copper phosphates are also mentioned, especially for use as food additives. The most effective combinations are indicated to be those in which both metals are present in the ionizable and nonionizable forms. The only testing reported in the French patent was odor reduction of sewage sludge, and odor reduction in mice by measurement of the odor of the entire mouse, fecal odor and urine odor.\nGerman Patent No. 1,083,503 describes the deodorizing action of metal complexes of 1,3-diketones. Listed metals include copper, nickel, cobalt, calcium, zirconium, zinc, tin, aluminum, cadmium, cerium, beryllium, magnesium, and mercury. Their use in various ways, including topical application to the axilla, is discussed. These compounds are theorized to function by interfering with the metabolism of the odor-producing microorganisms.\nBritish Patent Application No. 1,581,586 describes a sanitary foot wear article which includes a composition comprising copper, silver, or a copper-silver alloy powder dispersed in and held by water-insoluble resin binder. The British patent application indicates a belief that the metal reacts with substances secreted by the foot to produce metal salts which act as astringents and also act to prevent the growth of microorganisms.\nDespite the above, attempts have continued to find compositions which are effective against axillary malodor even after it has been formed."}
-{"text": "There are many and diverse types of apparatus which function to move an object in various directions. Among these apparatus are included power operated implements such as buckets, booms and bulldozer blades which are usually mounted on a vehicle and manipulated by a plurality of individual controls. The operator must give a great deal of attention to selectively moving the various levers, buttons, or the like into their proper modes of operation in order to correspondingly maneuver the various hydraulic motors or jacks which control the complex movement of the implement.\nFor example, while two separate levers can be utilized for individually operating a pair of valves in a hydraulic system and thereby directing fluid selectively to the opposite ends of a pair of hydraulic jacks, there is frequently the need to simultaneously move the jacks in the same direction or in opposite directions at the same rate. This is difficult to accomplish when the operator is also maneuvering the vehicle or manipulating a separate control element since two hands are usually required. On the other hand, if an attempt is made to mechanically connect each of the levers with both of the control valves so that the two jacks can be simultaneously extended or retracted with one lever or can be simultaneously moved in opposite directions with the other, then movement of one lever would disadvantageously move the other. In addition, it is usually desirable that both control levers be automatically returned to a neutral or jack-holding position upon manually releasing them, and yet when operating an individual lever the effort which is required should be minimal. It is thus apparent that a complex coupling problem is involved when connecting both levers to both valves.\nWhile it would appear to be desirable to utilize a single control lever to position a pair of hydraulic jacks through a suitable control valve arrangement, it is not always easy to relate the required hand movement with the desired implement movement. Furthermore, operator movement of a universally pivotable single control lever can easily lead to unequal rates of movement of the jacks to the point of requiring correctional manipulation of the lever at a crucial operating time.\nReference is herein made to U.S. Pat. No. 3,705,631 issued Dec. 12, 1972 to D. H. Seaberg, and U.S. Pat. No. 3,795,280 issued Mar. 5, 1974, the latter of which is assigned to the assignee of the present invention, which relate to control arrangements and associated hydraulic control circuits for selectively actuating separate hydraulic jacks for lifting, tilting and angling a bulldozer blade on a tractor. While such controls have been well received by the industry, it takes a special degree of dexterity and alertness to operate them. Particularly, it is difficult to correctly manipulate an electrical switch disposed on or adjacent to the single control lever. As a result, the operator will occasionally move the switch in the wrong direction or hold it engaged an excessively long period. As a consequence, one of the hydraulic jacks might well be at the end of its stroke so that the additional fluid flow being directed thereto by the inadvertent engagement of the switch must be exhausted at the relatively high pressure setting of the system relief valve. As is well known in the art, any prolonged exhausting of fluid against the system relief valve pressure is undesirable from many standpoints.\nAnother problem is that these lever control arrangements cannot always be disposed immediately adjacent to the control valves. As a result they must cooperate with various hydraulic or electrical piloting devices to remotely actuate the control valves and to thereby affect operating economies for certain member placement advantages on the vehicle."}
-{"text": "The safety of the soldiers in the field is high on the priorities of weapons design. Not only are soldiers constantly exposed to enemy fire and a hostile environment during battlefield engagement, they are also exposed to the firing hazards of their own weapons and launch equipment. Often times, the launch equipment and munitions may be worn out or become unsafe due to intense use, operator error, improper maintenance or poor weather conditions, resulting in an increase of potential hazards to the soldiers. It is therefore an important goal to minimize the potential hazards to soldiers operating weapons in the field despite such real world adverse conditions.\nIn one conventional design of an artillery gun, the recoil buffering mechanism for the artillery gun comprises a breech assembly and a barrel, wherein the barrel and breech do not provide any special safety design for the operator standing behind the breech to fire the artillery gun. This configuration leaves the operator fully exposed to the dangers of explosives and hot combustion gases leaking from the gun in case the breech accidentally opens during the gun firing.\nMore specifically, when the breech is in an open position, the munition is loaded axially into the firing chamber. To perform this operation, the operator positions himself or herself at a distance from the breech. Next, the breech is closed in preparation for firing. Since the firing event is accomplished almost instantaneously, the operator remains at the distal position, behind the breech, during the entire firing operation.\nA significant pressure rise results from firing the munition. It is necessary for the breech to remain safely closed during the firing in order to impart the maximum forward momentum to the projectile, and to prevent any of the explosives and hot combustion gases from leaking past the breech to cause harm to the operator along the leakage path.\nHowever, due to wear, debris or other unforeseen factors, the breech might not be fully closed prior to, or during the firing of the weapon, resulting in leakage of explosives and hot combustion gases from the barrel to a position distal to the breech where the operator is positioned. This exposure increases the hazard of the operator and could pose substantial danger to operator safety.\nIn one embodiment, of a safety and arming mechanism for a rifled gun, the mechanism is controlled by three projectile parameters. The first and second parameters are the axial and angular accelerations of the fired projectile, which move a setback ball to arm the mechanism. The third projectile parameter, i.e., angular velocity, is utilized to lock the setback ball in the armed position. As the projectile continues its flight, it only becomes armed when a spin actuated escapement mechanism is moved to a fully armed position. In addition, a command arm signal is required to release the arrangement such that the escapement mechanism is in a condition to complete its motion to the fully armed position.\nSuch weapon arming safety design is implemented in the munition only and not in the launch equipment. It does not explicitly address safety against explosive hazards or firing hazards at the point of firing in case a catastrophic failure occurs, as in the case of firing artillery rounds, and the breech suddenly becoming loosened from the closed position, leaking explosives and hot combustion gases behind the breech.\nAnother conventional safety-and-arming device is based on micro-electromechanical system (MEMS). Two independent mechanical locks are moved out of the way to allow the arming slider to remove a barrier in the explosive train to arm a fuze or close a switch for firing. The mechanical locks respond only to valid launch or deployment conditions. In addition, the mechanism does not explicitly address safety against explosive hazards at the point of firing in case a catastrophic failure occurs.\nIn yet another conventional device, a projectile is launched with on-board linear acceleration sensors to measure at least two accelerations, and the recorded time interval between the two accelerations would need to fall within a pre-determined range in order to arm the munition for detonation. This device assures the safety of arming the munition as long as the launched projectile achieves target values in flight parameters. When this goal is not achieved, the munition in flight would not be allowed to detonate. However, this device deals with the safety to arm the projectile after becoming airborne, and not with the safety of the weapon system during the firing process to protect the weapon operators.\nIn still another embodiment, a firearm safety locking mechanism prevents accidental or unauthorized use of the weapon. The safety locking mechanism is placed and operates in the firing chamber or in the barrel of the weapon. One of the goals of this mechanism is to prevent accidental use by an under-aged operator. However, such mechanism does not address the firing hazard reduction in case the firing chamber fails to hold the hot explosive gases in place inside the weapon.\nAlthough these conventional technologies have proven to be useful, the issue of safety at the point of firing has not been addressed, and it would be desirable to present additional improvements to further reduce firing hazard. What is needed is an artillery gun equipped with a breech having a mechanism to safeguard against premature firing of munitions before the breech is fully closed. The safety mechanism should prevent explosives and hot combustion gases from the primer and the charge from quickly leaking from the firing chamber past the breech and subjecting to harm any personnel in the path of leakage. The need for such a safety mechanism has heretofore remained unsatisfied."}
-{"text": "Compounds based on the ergoline ring system: ##STR1## have a suprising variety of pharmaceutical activities. For example, lysergic and isolysergic acid are D-8-carboxy-6-methyl-.DELTA..sup.9 -ergolines (9,10-didehydroergolines or 9-ergolenes.) The amides of lysergic acid have valuable and unique pharmacologic properties, and include the naturally-occurring peptide alkaloids; ergocornine, ergokryptine, ergonovine, ergocristine, ergosine, ergotamine, etc.; synthetic oxytocic alkaloids such as methergine; and the synthetic hallucinogen - lysergic acid diethylamide or LSD. Ergotamine, a 9-ergolene, with a \"peptide\" side chain, has been used in the treatment of migraine and recently, both ergocornine and 2-bromo-.alpha.-ergokryptine have been shown to be inhibitors of prolactin and of dimethylbenzanthracene (DMBA)-induced tumors in rats, according to Nagasawa and Meites, Proc. Soc. Exp'tl. Bio. Med. 135, 469 (1970) and to Heuson et al., Europ. J. Cancer, 353 (1970). (See also U.S. Pat. Nos. 3,752,888 L and 3,752,814).\nNon-peptide ergot derivatives, both naturally occurring and totally or partially synthetic, share these multiple pharmacological properties with the peptide derivatives. For example, D-6-methyl-8-cyanomethylergoline, was prepared by Semonsky and co-workers, Coll. Czech. Chem. Commun., 33, 577 (1968), and was found to be useful in preventing pregnancy in rats - Nature, 221, 666 (1969). (See also U.S. Pat. No. 3,732,231) - by interfering with the secretion of hypophysial leuteotropic hormone and the hypophysial gonadotropins or by inhibiting the secretion of prolactin. [See Seda et al., Reprod. Fert., 24, 263 (1971) and Mantle and Finn, id. 441)]. Semonsky and co-workers, Coll. Czech. Chem. Comm., 36, 220 (1971), have also prepared D-6-methyl-8-ergolinylacetamide, a compound which is stated to have anti-fertility and anti-lactating effects in rats. The 2-halo derivatives of D-6-methyl-8-cyanomethylergoline and of D-6-methyl-8-ergolinylacetamide have been prepared and found to be active prolactin inhibitors (M. J. Sweeney, J. A. Clemens, E. C. Kornfeld and G. A. Poore, 64 th Annual Meeting Amer. Assoc. Cancer Research, April, 1973 - See also U.S. Pat. No. 3,920,664).\nD-2-chloro-6-methyl-8-cyanomethylergoline (generic name, lergotrile), one of the aforementioned 2-halo derivatives, can be prepared by chlorinating D-6-methyl-8-cyanomethylergoline with such halogenating agents as N-chlorosuccinimide, N-chloroacetanilide, N-chlorophthalimide, N-chlorotetrachlorophthalimide, 1-chlorobenzotriazole, N-chloro-2,6-dichloro-4-nitroacetanilide, N-chloro-2,4,6-trichloroacetanilide, and sulfuryl chloride. Lergotrile can also be prepared by chlorinating any of the intermediate compounds used by Semonsky et al. (loc. cit.) for the preparation of D-6-methyl-8-cyanomethylergoline including methyl dihydrolysergate or D-6-methyl-8-hydroxymethylergoline, or with the novel intermediate of Example 6 of U.S. Pat. No. 3,920,664 such as, D-6-methyl-8-methylsulfonyoxylmethylergoline.\nIt is an object of this invention to provide a method of chlorinating any of the above ergoline derivatives to provide either lergotrile or a compound readily convertible thereto in higher yield than with any of the chlorinating agents heretofor provided by the art."}
-{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to automatic convergence of stereoscopic images based on disparity maps.\n2. Description of the Related Art\nIn human visual systems or stereoscopic camera systems, the point of intersection of the two eye axes or two camera axes is the convergence point. The distance from the convergence point to the eye or camera is the convergence distance. For human eyes, the convergence point can be at an arbitrary distance. For stereoscopic cameras, the convergence point may be, for example, at infinity (for a parallel camera configuration) or at a fixed distance (for a toe-in camera configuration).\nWhen a person looks at stereoscopic images on a stereoscopic display, the eyes naturally converge to the display screen. The distance from the display screen to the eyes is the natural convergence distance. However, to view the 3D effect correctly, the viewer's eyes adjust to have the same convergence distance as the camera. Such constant convergence distance adjustment can cause discomfort over time such as headaches or eye muscle pain."}
-{"text": "Field of the Invention\nThe present invention relates to a water output device, and in particular to a temporary stop water output device used in shower equipment.\nThe Prior Arts\nUsually, in using faucet water, sometimes it is required to stop using water temporarily, and then resume using water right after. However, presently, this is realized through a control switch not disposed on the water output equipment itself, and that could cause quite inconvenience.\nTherefore, presently, the design and performance of the water output device is not quite satisfactory, and it leaves much room for improvement."}
-{"text": "The present invention relates to a hydrogel device containing an active agent which is released by diffusion to an aqueous medium upon contact therewith at a controlled rate for a predetermined period of time, a method for preparing such a device and a method of use of such a device to release the active agent into the aqueous medium.\nVarious devices are known in the art for the sustained release of an active agent. For example, monolithic devices for sustained release of agent, wherein the agent is dispersed uniformly in a non-swellable homogeneous and imperforate polymer matrix where the agent dissolves in and permeates through the polymer itself are known. Microporous devices, in which the pores contain an active agent permeable liquid or gel medium such that the active agent preferentially dissolves in and permeates through the medium in the pores, are also known, as well as osmotic bursting devices wherein water is imbibed osmotically into active agent depots in a water permeable-active agent impermeable polymer matrix to rupture the depots serially. Characteristically, such devices release the active agent at a high initial rate which then drops off in the manner of a first order equation.\nThe release of an active agent from a uniformly dispersed agent in a non-swellable homogeneous and imperforate plasticized polymer matrix by active agent permeation can be classified as a classical Fickian release, since the release pattern generally follows Fick's law of diffusion; EQU M.sub.t =Kt.sup.1/2\nwhere M.sub.t is the amount active agent released, K is a constant and t is time.\nU.S. Pat. No. 3,923,939 describes the removal of a portion of active agent from the surface of the aforementioned prior art diffusion and osmotic bursting devices by washing to form a depleted layer of polymer matrix. However, the devices do not include hydrogels, and moreover still generally exhibit a Fickian release, albiet with a lower initial release as evidenced by the release profile of the diffusion type devices of Example 1A with those of 1B in FIGS. 3 and 4 of that patent. The only anomolous release pattern described and exemplified is that of an osmotic bursting device in which the bursting device was washed for 24 hours such that the initial release rate was decreased by 90% as described in Example 3B and FIG. 10 of U.S. Pat. No. 3,923,939.\nOne drawback of the drug diffusion devices of the type exemplified in Example 1B of the aforementioned patent is that after treatment of the surface, the drug will diffuse at a rate corresponding to the drug diffusion rate in the device matrix until the drug is again uniformly distributed throughout the diffusion device. As a result, such a drug soluble rubbery polymer system exhibits numerous drawbacks.\nActive agent releasing devices containing a swellable hydrogel matrix are also known in the art. In such devices, the active agent is uniformly distributed in the non-swollen, or glassy, hydrogel polymer. Upon contact with an aqueous environment, the dry hydrogel swells as water penetrates the glassy matrix. The boundary between the glassy phase and the swollen, or rubbery, phase of the hydrogel is known as the solvent front. As the aqueous solvent front moves inward from the surface to the center of the hydrogel device, the active agent in the rubbery swollen phase of the hydrogel dissolves and diffuses through the swollen phase into the external aqueous environment. The active agent is substantially incapable of diffusion to any significant extent in the glassy non-swollen phase of the hydrogel. The rate of active agent transport into an aqueous environment is dependent upon a number of factors including the rate of penetration of the aqueous solvent front, the shape of the device, the diffusivity of the active agent in the swollen phase, the amount of active agent loading in the hydrogel matrix, the distance between the solvent front and the surface of thw device, the rate of decrease of the surface area of the glassy phase during solvent penetration and the like. See, for example, Ping I. Lee, Polymer Comunications, Vol. 24, p. 45-47 (1983).\nCharacteristically, the release rate for such hydrogel devices, containing uniformly distributed active agent, generally conforms to the following equation EQU M.sub.t =Kt.sup.a\nwhere M.sub.t is the amount of active agent released, K is a constant, t is time and a is between about 0.8 and 0.5. As a result, such systems do not provide for a substantially zero-order release, since the rate of release continuously drops off from a high initial release rate."}
-{"text": "The present invention concerns that of a new and improved swing covering accessory. The accessory is used to cover an outdoor infant/toddler swing and is designed to be placed into the swing and around the outside of it, providing both entertainment and sanitary purposes for infants and toddlers using the swing."}
-{"text": "It is known from e.g. U.S. Pat. No. 4,025,575 that light oxygenates can be converted to the lower olefins ethylene and propylene using H-ZSM-5 zeolite as a catalyst.\nU.S. Pat. No. 3,911,041 describes that methanol and DME can be converted to a reaction product containing olefins using a zeolite catalyst containing at least 0.78 wt-% phosphorus incorporated with the crystal structure of the zeolite. The zeolite used in the process of U.S. Pat. No. 3,911,041 may be ZSM-5. U.S. Pat. No. 3,911,041 further teaches that the activity of the phosphorus-containing zeolite catalyst can be increased by depositing zinc (Zn) on the zeolite.\nU.S. Pat. No. 5,367,100 describes a method for the conversion of methanol or dimethyl ether to light olefins using a zeolite ZSM-5 based catalyst containing at least 0.7 wt-% phosphorus and at least 0.97 wt-% rare earth elements which are incorporated within the structure of the catalyst. The rare earth elements are preferably rich in lanthanum so that the content of lanthanum in the catalyst preferably is between 2.5 and 3.5 wt-%.\nU.S. Pat. No. 4,049,573 describes a number of different boron or magnesium-comprising catalyst compositions useful for converting monohydric alcohols and their ethers to a hydrocarbon mixture ring in C2-C3 olefins and mononuclear aromatics.\nThe oxygenate-to-olefin catalysts of the prior art have the disadvantage that they have a relatively low selectivity for C2-C3 olefins and/or become quickly deactivated by coke deposition on the catalyst surface."}
-{"text": "The present invention relates to integrated circuit (IC) design, and more specifically, to a method and system for clock tree construction, an IC and a fabrication method thereof.\nA typical IC may comprise a large amount of logic elements and other circuits for implementing IC logic functionality. Further, an IC chip may comprise a clock tree (i.e., a clock signal distribution network) for distributing a clock signal received at an input to all clock sinks that are \u201cclocked\u201d by the clock signal. A clock tree may comprise wires, buffers etc, to distribute the \u201cclock signal\u201d that controls the timing and operation of logical elements and other circuits of the IC. Clock sinks (or sinks) refer to logic elements or other circuits such as registers (flip-flops), RAM and latches, controlled by a clock signal, such that they add capacitance to the clock tree. Those sinks can change their states in response to clock signal pulses, and the IC synchronizes state changes in various sinks in a clock domain by clocking them with the same clock signal.\nClock skew is a significant aspect in assessing clock tree performance and quality. Clock skew generally refers to the difference (delay) between arrival times at any two clock sinks of a clock signal from an external clock source. Due to different path lengths of various branches to the respective clock sinks in a clock tree, there may be some clock skews between those various clock sinks. Further, in order to deliver the clock signal to every region of an IC, clock cells (buffers, for example) are generally inserted in the clock tree to amplify and/or retransmit the clock signal. However, because each clock cell has an intrinsic delay, it may cause a certain clock skew also. Thus, controlling or restricting the level number of buffers in a clock tree is one way for improving clock tree performance and IC design quality. Theoretically, smaller clock skews can be obtained, if there are less but the same levels of buffers contained in each branch leading to various clock sinks in a clock tree. However, the above assumption cannot be satisfied in practical IC design in many situations. With the technical evolution of digital IC design, a common path method is becoming more important for improving clock tree's skew and timing. A common path generally refers to a path consisted of buffers that are shared by multiple sinks in a clock tree. The longer the common path is, the smaller the clock skew of a clock signal arriving at sinks is. Traditional techniques employ a method to maximize the common path, that is, to allow sinks in a clock tree to share buffers at various levels as much as possible. In principle, the more buffers shared in a clock tree, the longer the common path will be. As a result, the performance of the clock tree may be optimized, and the quality of the IC designed may be improved. Other means for clock tree optimization include, for example, utilizing high performance clock cells (elements) capable of reducing clock skew, and the like."}
-{"text": "The present invention relates to an external cavity laser on silicon, and more specifically, to temperature insensitive external cavity lasers on silicon.\nAn optical cavity or optical resonator is an arrangement of mirrors that forms a standing wave cavity resonator for light waves. Optical cavities are a major component of lasers, surrounding the gain medium and providing feedback of the laser light. They are also used in optical parametric oscillators and some interferometers. Light confined in the cavity reflects multiple times, producing standing waves for certain resonance frequencies. The standing wave patterns produced are called modes. Longitudinal modes differ only in frequency while transverse modes differ for different frequencies and have different intensity patterns across the cross section of the beam.\nDifferent resonator types are distinguished by the focal lengths of the two mirrors and the distance between them. Flat mirrors are not often used because of the difficulty of aligning them to the needed precision. The geometry (resonator type) must be chosen so that the beam remains stable, which means that the size of the beam does not continually grow with multiple reflections. Resonator types are also designed to meet other criteria such as minimum beam waist or having no focal point inside the cavity. Optical cavities are designed to have a large Q factor, which means that the light beam will reflect a very large number of times with little attenuation. Therefore, the frequency line width of the beam is very small compared to the frequency of the laser.\nLight confined in a resonator will reflect multiple times from the mirrors, and due to the effects of interference, only certain patterns and frequencies of radiation will be sustained by the resonator, with the others being suppressed by destructive interference. In general, radiation patterns which are reproduced on every round-trip of the light through the resonator are the most stable, and these are the eigenmodes, known as the modes, of the resonator.\nResonator modes can be divided into two types: longitudinal modes, which differ in frequency from each other; and transverse modes, which may differ in both frequency and the intensity pattern of the light. The basic or fundamental transverse mode of a resonator is a Gaussian beam.\nThe most common types of optical cavities consist of two facing plane (flat) or spherical mirrors. The simplest of these is the plane-parallel or Fabry-P\u00e9rot cavity, consisting of two opposing flat mirrors. Plane-parallel resonators are therefore commonly used in microchip lasers, microcavity lasers, and semiconductor lasers. In these cases, rather than using separate mirrors, a reflective optical coating may be directly applied to the laser medium itself."}
-{"text": "1. Field of the Invention\nThe present invention relates to an apparatus to control the ratio of air to fuel of the air-fuel mixture which is supplied to an internal combustion engine, and more particularly to an apparatus provided for controlling the amount of air supplied to the main system fuel passage and the slow system fuel passage of a carburetor for supplying an air-fuel mixture to an internal combustion engine by means of a single proportional control solenoid valve.\n2. Description of the Prior Art\nIt is known to incorporate a catalytic converter have ternary catalytically active substances into the exhaust system of the internal combustion engine of an automobile in order to simultaneously reduce the injurious components of the exhaust gas, such as hydrocarbons (HC), carbon monoxide (CO) and nitrogen oxides (NOx). It is required to supply an air-fuel mixture of an air-fuel ratio corresponding to the stoichiometric air-fuel ratio into the cylinders of an internal combustion engine, since the maximum cleaning efficiency of the ternary catalytically active substances is attained when the exhaust gas is produced by the combustion of an air-fuel mixture of the stoichiometric air-fuel ratio. An air-fuel ratio controller to meet such a requirement is proposed, which includes a main system air-bleeding passage connected to the main system fuel passage connecting to the main jet of a carburetor, a slow system air-bleeding passage connected to the slow system fuel passage connecting to the slow system fuel supply port of the carburetor, and a main system proportional control solenoid valve for controlling the main air bleeder and a slow system proportional control solenoid valve for controlling the slow air bleeder, which are disposed within the main system air-bleeding passage and the slow system air-bleeding passage, respectively, and are adapted to be controlled by a control signal provided by converting the output signal of an oxygen sensor disposed within the exhaust passage of the engine by means of an electronic control unit. It is proposed to control the rate of air to be supplied to the main system fuel passage and to the slow system fuel passage through the main system air-bleeding passage and through the slow system air-bleeding passage, respectively, with an air-fuel ratio controller as described above, so that the air-fuel ratio of the air-fuel mixture which is supplied into the cylinders of an engine is controlled so as to be close to the stoichiometric air-fuel ratio.\nHowever, heretofore known air-fuel ratio controllers of this type for carburetor are provided with individual proportional control solenoid valves in the main system air-bleeding passage and the slow system air-bleeding passage, respectively, and are adapted to control both proportional control solenoid valves simultaneously with the control signal of a single system provided by the control signal producing circuit of an electronic control unit. Accordingly, the electronic control unit is required to supply an electric current simultaneously to the main system and the slow system proportional control solenoid valves. Therefore, the electronic control unit is necessary to be capable of supplying an electric current twice as much as that to be supplied to a single solenoid. Thus such conventional air-fuel ratio controllers have a disadvantage that the control signal producing circuit of the electronic control unit must consist of elements which have superior current capacities."}
-{"text": "1. Field of the Invention\nThe present invention is directed to process controllers. More specifically, the present invention is directed to a process controller for controlling a process in response to the duration of a controller disabling power outage."}
-{"text": "The present invention relates to a honeycomb structure with a sensor insertion hole."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to information retrieval systems and, more particularly, to natural language information retrieval systems.\n2. Related Art\nAs information systems become increasingly interconnected through intranets and internets, the main problem with the search for information has shifted from determining whether the requisite information exists to determining how to locate such information.\nTo explore the Internet or some other large database, well-known browsers and search engines are available. Unfortunately, currently existing search engines generally require the use of expressions and search methods with which the user has to be familiar, and generally require the user to enter keywords assumed to be related to the information that the user seeks. Moreover, a typical search returns a vast amount of information, much of it being irrelevant to the user. The user is then required to find the few relevant documents from the search results. Furthermore, it is not uncommon to not locate interesting documents because the user did not use keywords that corresponded to the words or word forms as written in those documents. Various conventional techniques have been developed to improve the recall of these and other information retrieval systems. However, these techniques have many drawbacks that limit their effectiveness in improving the ability of information retrieval systems to identify all available information related to a desired search topic.\nA primary drawback is that the conventional method for expressing the content of a text is a single word extraction. Conventional information retrieval methods rely on word stemming or rooting, skip word filtering, and proximity measures. Typically, conventional systems stem or root the words that occur in the text and subsequently filter out all stems or roots that appear in a predetermined skip list. The skip list contains words that have little or no predictive value. Such words include function words, such as articles, pronouns, prepositions and other frequently used words. The end result of these conventional methods is a keyword list containing a list of single non-trivial words occurring in the document, optionally ordered by their frequency in the text. The keywords in the keyword list are in their stemmed or rooted form, and accompanied by their offset values or similar location markers.\nFor computer applications that require an intelligent representation of the content of a text, these conventional methods are inadequate. For example, meaningful units of content generally consist of more than one single word as provided by the conventional information retrieval systems. For example, in a keyword list containing, among others, \"Amsterdam\", \"Rotterdam\", \"Marathon\", and \"Airport\", the informational content of the text, the Rotterdam Marathon and Amsterdam Airport, is lost. Likewise, in a query \"are the Antwerp Yellow Pages in the Web yet?\", the crucial phrase evidently is \"Antwerp Yellow Pages\". This phrase needs to be parsed and processed in a way that retains the informational content of the query. Specifically, only those documents that literally match on \"Yellow Pages\" and not just any occurrence of either \"page\" and \"yellow\" separately, as well as on \"Antwerp\" in a premodifying or postmodifying position, rather than just in any location in the text, should be retrieved from the searched database. However, conventional information retrieval systems typically yield the query \"Antwerp or Yellow or Page\" which may retrieve, among others, documents on Flemisch paper factories. Furthermore, the single keywords \"yellow\" and \"page\" fail to express the notion that emerges from their combination and preservation of the plural form; that is, \"Yellow Pages\".\nIn addition, the single keyword lists used in conventional information retrieval systems do not merge expressions that are different in form but share the same reference. A method that ignores synonyms, hyponymys, name variants, frequent misspellings, and other semantic relations, fails to give a proper representation of the document's content. For example, a text dealing with the wife of the current president of the United States may contain any number of references to that person, ranging from \"the President's wife\" and \"The First Lady\" to \"Mrs. Clinton\" and \"Hilary.\" Conventional information retrieval systems ignore these synonymous expressions, failing to give a proper representation of a text's contents.\nAnother drawback to conventional systems is that their mechanical application of a skip word list ignores the content representation of the text. For example, the individual word \"page\" may be a skip word, but in combinations like \"Yellow Pages\" or even \"The Sports Pages,\" it should be preserved. Likewise, skipping \"first\" from \"First Lady\" leads to a loss of the essence of the expression.\nWhat is needed, therefore, is an apparatus and method for efficiently retrieving information that accurately represents the content of both the text being searched, and the user's query, in such a way that the two can be more effectively matched."}
-{"text": "Prior to the conception and development of the present invention, as is generally well known in the prior art, light bulbs are used in practically every room of a home or office and must be replaced every so often. Climbing up a ladder or balancing on a chair to change burned out bulbs in a high ceiling fixture can be quite difficult and dangerous for anyone, even possibly resulting in injury to those involved."}
-{"text": "The disclosure generally relates to small containers or boxes, such as those used for displaying or packaging jewelry and other items. Typically, jewelry, watches, and other such items are displayed on a store rack, which may be placed on a counter. Exemplary racks used for display may include one or more flat display faces that may be mounted on a rotatable base to enable a customer to quickly look through the items on display.\nEach display face of a typical display rack may include a plurality of regularly spaced hooks or openings, each of which is configured to support a particular item for display. Items are usually placed within boxes that are either open or transparent to allow the customers to view the contents. Open boxes in particular are used to display jewelry and other items that have visual appeal. To support the boxes on the rack, disposable attachments to the boxes, such as removable plastic hooks, have been used in the past. For example, a removable plastic hook that attaches into an opening of the box may be used to hang the box from an opening in the display rack. When an item in its box is selected by a customer, a sales associate may remove the box from the display, remove and dispose of the hook, and retrieve a lid or other closure of the box such that the item may be purchased in its packaging, which may be gift wrapped.\nThis type of known product display and packaging arrangement presents certain disadvantages. For one, the cost of merchandise packaging is increased by the use of the disposable plastic clips. Further, waste is created by the disposable clips used to mount the open boxes to the display. Even further, the lids for the open display boxes are stored near the display for use when items are selected for purchase, thus taking up space and time for the sales associate to retrieve them."}
-{"text": "In certain types of building construction, e.g., retrofits, cinder block, or the like, it may not be possible or practicable to run electrical cabling (such as high voltage power lines and low voltage data lines) through the building's walls. In such cases, modular raceway or conduit assemblies are often used to house and route cabling along a wall or other surface. A typical raceway assembly includes a linear or elongate housing having at least one interior passageway that accommodates a length of electrical cabling. (Such a passageway is referred to herein as a \u201cwireway.\u201d) The housing is attached to a wall, and then the cabling is disposed in the interior of the housing. To cover a given span, multiple segments of housing are deployed in an end-on-end manner.\nIn a typical raceway or conduit installation, sections of conduit must be cut precisely to match the wall span to be crossed by the raceway. Such cutting operations, e.g., through sheet metal, are time consuming for installers. Additionally, because the sections are custom cut, it is difficult to adjust or compensate for variations in spacing that stem from adding junction boxes, retrofit outlets, and the like."}
-{"text": "To improve the mechanical and physical properties of castings, it is well known in the art to use inserts which can locally improve the strength, wear resistance or other characteristics of the casting. The inserts are placed within the mold prior to pouring the molten casting material and, upon cooling of the molten casting material, the inserts form an integral part of the finished cast product. A common example of such an application is in the engine block of an internal combustion engine. In particular, aluminum alloy engine blocks, typically an aluminum-silicon alloy, often make use of cast iron inserts, or liners, which are more durable than the cast aluminum walls of the engine block, particularly when aluminum pistons are used. A process well known to those skilled in the art is referred to as the \"Al-fin\" process, which entails dipping a cast iron insert in molten aluminum prior to placing the insert in a mold and casting a molten aluminum around the insert.\nHowever, a disadvantage with the use of cast iron liners is the significant additional weight which is incurred. In addition, cast iron does not conduct heat away from the cylinder as well as aluminum, which raises the temperature of the cylinder and imposes higher temperature-related stresses and wear on the engine's internal components. Another disadvantage with using iron is that there is a mismatch between coefficients of thermal expansion between iron and aluminum and its alloys, which can cause debonding of the insert.\nAs a result, it is generally preferable to provide inserts which are lower in weight while also providing better heat transfer capability and a more closely matched coefficient of thermal expansion. Naturally, aluminum-base alloys and composites are generally suitable in terms of weight, heat transfer and thermal expansion, for use with aluminum castings. Unfortunately, aluminum-base inserts do not metallurgically bond well to casting materials because the insert forms an aluminum oxide layer at the insert's surface. The presence of oxides produces a weak bond because of the inability of the molten casting material to wet the insert's surface.\nTo overcome this problem, one approach known in the art has been to form an insert which can be penetrated by the casting material under high pressure to form a mechanical/metallurgical bond between the casting material and the insert. One such approach uses alumina and carbon fibers which have been highly compressed to form a cylindrical insert. An aluminum casting material is then pressurized sufficiently during the casting process to penetrate the fiber inserts without structurally damaging them. While durability is improved, manufacturing costs are significantly higher than that of iron liners in aluminum blocks.\nAn approach for promoting a metallurgical bond between the insert and the casting material is taught by U.S. Pat. No. 4,687,043 to Weiss et al. Weiss et al. provide an aluminum composite insert whose outer surface is covered with an aluminum alloy. The insert is then coated with a molten solder alloy. The molten solder alloy is selected to have a melting temperature which is below the melting temperature of the insert's aluminum alloy cover layer. The insert is dipped into the molten solder alloy to separate from the cover layer the oxides already present and to prevent the formation of new oxides thereon. In addition, Weiss et al. teach that the casting material must be at a temperature which is higher than the melting temperatures of the insert's cover layer and the solder alloy. This enables the casting material to flush the molten solder alloy from the cover layer during the casting process to expose the cover layer to the casting material, allowing the casting material and the cover layer to form an oxide-free metallurgical bond. The molten solder alloy is intended to be mixed with the casting material and not remain on the surface of the insert.\nIn addition, Weiss et al. teach a zinc solder alloy which contains about 10 to 30 weight percent tin and about 5 to 25 weight percent cadmium to reduce the melting temperature of the solder alloy below the temperature of the casting mold. A disadvantage to the use of tin for the purposes taught by Weiss et al. is that tin embrittles the solder alloy and the interface between the insert and the casting, while also reducing their corrosion resistance. The presence of tin also slows down the age hardening by reducing the Guinier-Preston (GP) zones formation rate, potentially causing a weak interface between the insert and the casting. Similar to tin, cadmium reduces corrosion resistance. But most importantly, cadmium poses an environmental concern in that it is highly toxic. As a result, the use of cadmium is always avoided where possible, and sometimes prohibited.\nThus, it would be desirable to provide a method of promoting a metallurgical bond between a strong, wear-resistant insert and an aluminum alloy casting material which would be economical for use in mass production, while also avoiding the concerns for embrittlement and toxicity."}
-{"text": "Automatically identifying the locations of objects and their parts in video is important for many tasks. For example, in the case of human body parts, automatically identifying the locations of human body parts is important for tasks such as automated action recognition, human pose estimation, etc. Body parsing is a term used to describe the computerized localization of individual body parts in video. Current methods for body parsing in video estimate only part locations such as head, legs, arms, etc. See e.g., \u201cStrike a Pose: Tracking People by Finding Stylized Poses,\u201d Ramanan et al., Computer Vision and Pattern Recognition (CVPR), San Diego, Calif., June 2005; and \u201cPictorial Structures for Object Recognition,\u201d Felzenszwalb et al., International Journal of Computer Vision (IJCV), January 2005.\nMost previous methods in fact only perform syntactic object parsing, i.e., they only estimate the localization of object parts (e.g., arms, legs, face, etc.) without efficiently estimating semantic attributes associated with the object parts.\nIn view of the foregoing, there is a need for a method and system for effectively identifying semantic attributes of objects from images."}
-{"text": "About one third of one's lifetime is spent sleeping and a not insignificant effort has been spent attempting to achieve comfort during this time period. Many inventors have attempted to approximate the comfort that one experiences in floating in a body of water such as the Great Salt Lake in Utah. It is believed that the California Indians spent the better part of cold winter days semi floating in pools of hot mud fed by the natural warm springs. The so called \"water bed\" has become popular, especially in the State of California. Water beds have taken two approaches; those that provide a \"floatation\" or \"buoyancy\" effect, and those that result in a \"hammock\" effect. The first efforts known to Applicant to provide a \"floatation\" effect is disclosed in White, U.S. Pat. No. 184,487, Nov. 21, 1876. White disclosed a plurality of water tight sacks instead of a \"single rectangular oblong air and water tight sack\". The White bed consisted of a rigid perimeter and the entire bed was filled with the water sacks. The rigid border probably provided a \"floatation\" effect. On June 15, 1971, Hall received U.S. Pat. No. 3,585,356 for a water bed consisting of a single flexible substantially inelastic bladder contained by a perimeter rigid frame. Hall, appears to be the first to fully explain the \"floatation\" effect which results from the use of a frame which provides lateral support to the flexible inelastic bladder. Hall filled the entire bed with water resulting in an extremely heavy piece of furniture. The entire body of the person was supported on the inelastic water bladder.\nThe \"hammock\" effect water bed is discussed in the Labianco patent U.S. Pat. No. 3,840,921. Labianco taught that if a person lies on a pillow type water bag which is not supported on its sides, the top of the water bag is placed in tension and \"it will create conventional bed pressure points, and will not conform and adjust properly to the different weight proportions of the upper and lower torso of the user.\" In effect, the body is supported in a \"hammock\" type attitude.\nBoth the laterally supported \"floatation\" beds and the \"hammock\" type non-laterally supported beds have the problem of excessive weight. More recently a great deal of activity has centered on the problem of reducing weight by providing air or foam rubber or plastic perimeters. Examples of such reduced weight beds are taught by Tobinick in his U.S. Pat. No. 3,702,484, U.S. Pat. No. 3,789,442 and 3,815,165; and Tinnel U.S. Pat. No. 4,015,299. All of these patents lack sufficient structure to provide the necessary lateral support to the flexible water container to result in a \"floatation\" type support. All of these beds which are combinations of unconnected foam materials, water containers and support sheathes provide a \"hammock\" type of support.\nReduction of mattress weight by replacing portions of the water area with lightweight foam materials has created a comfort problem in that a person lying on beds constructed in the manner of Tobinick, U.S. Pat. No. 3,789,442 and Tinnel, U.S. Pat.No. 4,015,299 can feel the edge of the foam material where it borders the cavity holding the flexible water bladder. This problem can be overcome, but not with foam mattresses which have cavities which do not provide sufficient lateral support. Hall, supra, taught that lateral support about the perimeter of the water bladder was essential. Without lateral support, a body resting therein sinks to an undesirable depth in the water. If the legs are resting on the foam at an elevation much higher than the heavier portions of the body such as the hips or shoulders, it gives the person an uncomfortable feeling and in addition he will feel the boundry between the foam and the cavity holding the water bag. This disparity in vertical displacement is even accentuated in Tobinick, U.S. Pat. No. 3,789,442 because his cavity walls are slanted.\nFinally, the prior art mattresses having foam material areas and water supported areas failed to recognize the great difference in evenness of body weight and the pressure points which occur at the hips and shoulders. Under these pressure points, foam material is more compressed than it is under the arms or legs. This unevenness in support can cause discomfort over protracted periods of time unless the sleeper shifts his weight. This unevenness of support is especially apparent in Tobinick U.S. Pat. No. 3,789,442 in which the hips and legs are supported on the water bag and the shoulders and arms are supported on foam material. While the Tobinick hospital bed had the advantage that the upper portion can be raised from the horizontal without disturbing the water cavity, the level of comfort in the hips and the shoulders is quite different."}
-{"text": "The invention described and claimed herein relates to the field of detecting rotation of a structure about an axis. In particular, the claimed invention relates to the field of monolithically-integrated gyroscopes."}
-{"text": "As a result of increased environmental awareness, water conservation has become an important issue especially in the Western part of United States, densely populated areas elsewhere within the United States and in other countries or regions that have arid climates.\nWith regard to domestic water consumption in this country, the ultimate use is generally for sanitary, culinary, drinking, washing or bathing purposes. It has been found that careless or willful waste occurs during the performance of certain of these tasks such as washing dishes, preparing food, brushing teeth, etc. wherein the user permits the water to flow continuously from the faucet rather than to temporarily shut-off the flow, when water is not needed, during these procedures. The primary reason for not interrupting the water flow is generally because the user's hands are occupied or because of the effort required to again achieve the same water temperature and/or pressure.\nBy way of example, water consumption during the process of washing and rinsing a typical load of dishes for four place settings, consumes 7.5 gallons of water or approximately one cubic foot of water. During the wash/rinse period, the water is idling for about fifty (50%) percent of the time, thus wasting 0.5 cubic feet.\nAside from the environmental issues, another motivation for reducing consumption is economic in nature in that there is a tendency for the rates charged for metered water to generally increase rather than decrease. Furthermore, the conservation of hot water minimizes the fuel needed to heat the water and thus achieves additional cost savings.\nAn apparatus to control a water faucet valve without requiring hands-on operation is shown in U.S. Pat. No. 2,270,239. A shortcoming of that device is that the moveable components of the mechanical linkage are subjected to friction and wear resulting in slack and inefficiency. A further problem with that device is that the installation requires a plumber or other skilled workman and it is not particularly adapted as a retro-fit assembly.\nAnother remote control faucet valve device, as shown in U.S. Pat. No. 4,052,035, utilizes a flexible conduit for communicating between a foot control member and a valve member. A disadvantage of that arrangement is that the conduit is draped over the sink and the kitchen cabinet when in use and thus presents a physical impediment or hinderance which interferes with access to the sink and underlying cabinet. The conduit can also further become entangled with kitchen appliances. Additionally, the hydraulic system has a restricted range of control functions.\nWith regard to wireless control systems utilizing a radio link, these devices have been applied to the operation of motorized toy vehicles and model airplanes. They have also been applied to remote keyless entry systems for locking or unlocking the doors and trunk lid of automobiles. The utilization of radio controls for home automation has been rather limited with the most common usage being for electronic garage door openers.\nFluid handling apparatus employing radio signal control has been applied to irrigation systems as discussed in U.S. Pat. Nos. 3,726,477 and 4,838,310. The systems of the aforementioned patents, however, are not adapted for domestic use and do not include an electronic valve structure or a tactile foot-operated sending unit as in this invention."}
-{"text": "The invention generally relates to automotive safety restraint systems and more particularly to a pretensioning system capable of simultaneously removing the slack in a plurality of seat belts, secured about respective occupants, utilizing a single pretensioning or belt-tightening device.\nPretensioners or belt tighteners, as they are also known in the art, are used to tighten a portion of a seat belt system, thereby removing slack about the occupant, during an accident. This action is done prior to the time the occupant begins to move forward as a result of the deceleration of the vehicle.\nAs a typical accident progresses, the occupant will remain fixed to the seat and begin to move forward at about 20 milliseconds after the initiation of the crash. Belt tighteners must react and eliminate the slack in the seat belt system, about the occupant, prior to the occupant's moving forward. Typically the pretensioning system is designed to react quickly and eliminate the slack in about 10 milliseconds upon receipt of a control signal, which is generated by one or more crash sensors and associated electronic control systems. There exists in the art two basic types of seat belt pretensioning systems. The first is associated with the spool of a retractor and the shoulder belt of a three-point seat belt system. During the accident the spool is caused, by various means, to reverse-wind or retract thereby removing the slack, typically in the shoulder belt portion of the seat belt system, to thereby hold the occupant in place. The other type of pretensioning system is called a buckle pretensioner. The buckle pretensioner is typically connected by a length of seat belt webbing or by a cable to the seat belt buckle. These buckle pretensioning systems (as well as retractor systems) include a pyrotechnic element which is fired, causing the buckle to be pulled downwardly, thereby reducing slack about the occupant. The prior art generally shows one pyrotechnic element associated with each shoulder belt (for the retractor pretensioner) or with each buckle (for a buckle pretensioner).\nIt is an object of the present invention to provide a single pretensioning device that is capable of eliminating the slack in a plurality of seat belt systems.\nAccordingly the invention comprises a system for eliminating residual seat belt slack about a plurality of occupants of a vehicle, the system comprising: a plurality of seat belt assemblies securable about each respective occupant, each seat belt assembly including a lap belt, a tongue, and a buckle. One end of the lap belt is secured to a support member and an opposite end of the lap belt is operatively joined to the tongue. The buckle matingly locks with the tongue and the buckle is secured to one end of a connecting member, the connecting member having an opposite second end. A belt tightening means, secured to the second ends of each connecting member, generally simultaneously tightens each seat belt assembly about a respective occupant. Each belt assembly may also include a shoulder belt.\nMany other objects and purposes of the invention will be clear from the following detailed description of the drawings."}
-{"text": "The present invention relates to new and useful N-phenyloxazolidinone compounds and their preparations, and more particularly to N-phenyloxazolidinone compounds in which the phenyloxazolidinone moiety is linked to a variety of saturated, or partially saturated, 4-8 membered heterocycles containing oxygen, nitrogen, and sulfur through a carbon-carbon bond.\nThe compounds are useful antimicrobial agents, effective against a number of human and veterinary pathogens, including gram-positive aerobic bacteria such as multiply-resistant staphylococci and streptococci, as well as anaerobic organisms such as bacteroides and clostridia species, and acid-fast organisms such as Mycobacterium tuberculosis and Mycobacterium avium. The compounds are particularly useful because they are effective against the latter organisms which are known to be responsible for infection in persons with AIDS."}
-{"text": "The present invention relates to novel telomers formed by the reaction of chlorotrifluoroethylene, herinafter designated as \"CTFE\", with fluoroxytrifluoromethane, hereinafter designated as \"FTM\". These telomers have properties which distinguish them from other CTFE telomers and make them commercially attractive products. Such properties include superior solvent characteristics which are useful in formulating non-flammable hydraulic fluids.\nVarious methods of preparing telomers from CTFE are known in the prior art and have been practiced commercially for many years.\nAn article by William T. Miller, Jr. et al in Industrial and Engineering Chemistry, pages 333-337 (1947), entitled \"Low Polymers of Chlorotrifluoroethylene\", describes a process for producing low molecular weight polymers of CTFE by carrying out the polymerization in a solution of chloroform using benzoyl peroxides as polymerization promoters. Other solvents disclosed in the reference as being useful for this purpose include carbon tetrachloride and tetrachloroethylene. The solution is heated in a pressure vessel for 13/4 hours at 100.degree. C., and the unreacted CTFE monomer and chloroform are removed by distillation, leaving a \"crude\" telomer of general formula CHCl.sub.2 (CF.sub.2 CClF).sub.n Cl, which can further heated and distilled to yield products ranging from a light oil to a semisolid wax or grease. In this formula, n is the chain length (the number of repeating units in the telomer chain), and is in the range of 1 to 20.\nA more recent development in this field is described in a series of articles by Y. Pietrasanta et al. entitled \"Telomerization by Redox Catalysis\" appearing in the European Polymer Journal, Vol. 12 (1976). This technology involves the reaction of a chlorinated telogen, such as carbon tetrachloride, with CTFE in the presence of benzoin and a suitable redox catalyst, such as ferric chloride hexahydrate (FeCl.sub.3.6H.sub.2 O). The telomerization reaction is suitably carried out in acetonitrile which is a common solvent for the reactants and catalysts. The telomerization reaction can be illustrated as follows, wherein n is as defined above: ##STR1##\nThe telomers produced according to the above-described process have end groups containing chlorine atoms. When used in a particular application, especially high temperature or corrosive applications, the chlorine atoms can be hydrolyzed, which may result in cleavage of the telomer and a loss of physical properties of the fluid. In order to prevent such a condition, the end groups are stabilized by fluorinating the telomer to replace chlorine atoms with fluorine. The fluorine atoms form a stronger bond with carbon and are less prone to cleavage. This results in a fluid which has superior performance over a wider range of operating conditions.\nFluorination of the telomer is accomplished by reaction with a suitable fluorinating agent, such as chlorine trifluoride or hydrogen fluoride. However, fluorination involves an additional process step which increases processing costs. It would be desirable to reduce the amount of fluorination required or eliminate this procedure in its entirey. It would also be desirable to develop superior telomers having imroved physical characteristics, such as better solvent properties. These are the primary objectives of the present invention."}
-{"text": "Commercial fishermen catch and process tons of fish a day. Many fishermen have complex, highly-automated processing equipment which dress the fish and produce numerous different commercial products for commercial retailers, such as restaurant chains or grocery stores. Such equipment is highly specialized to perform a specific task, such as cutting the heads and tails off, filleting, deboning, and skiing, etc. One of the reasons such equipment must be tailored to specific tasks is the large variety of shapes and sizes of fish processed, including pollock, salmon, trout, sole, cod, etc.\nTo reduce overall processing costs and to improve quality, commercial fishing vessels include processing equipment and freezers located directly on the ships. As the fish are caught, they are immediately processed and frozen, substantially improving the quality of the final product. This processing equipment also allows the vessels to remain at sea for long periods of time without transporting the fish to a shore-based processing facility. Most of the processing equipment is large, complex, and expensive. It is difficult and expensive to remove the processing equipment from the vessel and replace it with new processing equipment.\nBecause of the expense associated with exchanging the processing equipment, including vessel downtime, most commercial processing ships are intended to catch and process specific types of fish. As an example, a vessel may be fitted with processing equipment designed to process generally oval fish having a relatively thick cross-section, such as pollock or salmon, or generally flat fish having a relatively thin cross-section, such as sole.\nDue to the natural migration of fish, and various fishing laws, individual types of fish are caught and processed during limited parts of the year. Therefore, vessels having equipment capable of processing only one type of fish sit idle during part of the year. If the processing equipment could be readily exchanged or adapted to process both oval fish, such as pollock or salmon, and flat fish, such as sole, the vessels could be operated over a larger part of the year. This would in turn reduce the amount of downtime and associated expense for the vessels.\nOne manufacturer of fish processing equipment is Baader GmbH & Co. KG, Lubeck, Germany (\"Baader\"). One of the fish processing machines produced by Baader is the Baader 182. The Baader 182 is designed to process pollock and salmon shaped fish and includes a conveyor having a plurality of fish holders mounted along the length of the conveyor. Each fish holder includes a recess configured to receive a pollock. While operating, a worker removes fish from a fish hopper located next to the conveyor and places an individual pollock in the recess of each fish holder. The conveyor subsequently moves the fish holders and pollock along the length of the conveyor into contact with rotating blades that cut both the head and tail of the pollock off. The body of the pollock is then carried to additional equipment that cuts the pollock into fillets and removes the bones and skin.\nThe Baader 182 works well with fish having an oval cross-section, such as pollock, but does not allow other types of fish to be processed. The Baader 182 is particularly unsuitable for processing fish with a relative flat cross-section, such as sole. The fish holders used on the Baader 182 are not capable of holding such fish during processing.\nIn order to adapt the Baader 182 to process fish having a relatively thin cross-section, the inventor previously applied for and received a patent on a device for holding flat fish during processing. U.S. Pat. No. 5,358,441 entitled \"Device for Holding Flat Fish During Processing\" issued to the present inventor on Oct. 25, 1994. The device for holding flat fish during processing allows the Baader 182 to process flat fish in addition to pollock. One embodiment of the device included a generally rectangular fish holder that replaced the fish holders on the Baader 182. The fish holders include two arms which extend over a portion of the fins and body of a flat fish and hold it in place on the fish holder. The arms include two recesses that have slanted rear walls that center the flat fish as it is placed within the fish holder. Two slots extend across the width of the fish holder to allow rotating blades to extend beneath the upper surface of the fish holder to ensure that the blades cut cleanly through the thickness of the fish.\nOne of the products that the device for holding flat fish during processing allowed the Baader 182 to be used for is the processing of \"kirimi.\" In order to make kirimi, the head and tall of a sole are removed, leaving the center portion of the sole. The center portion of the sole is then sold to consumers who fry, bake, or otherwise prepare the center portion for consumption. In the past, the center portion was sold complete with any internal organs remaining within the center portion. One of the internal organs remaining in the center portion after removing the head and tall is the kidney, commonly referred to as the \"blood spot.\" The blood spot is located along the backbone of the sole in the internal organ cavity. In the past, the blood spot has been either removed by the individual consumer during preparation or removed by hand during processing after the head and tail are removed.\nIn today's highly competitive fishing industry, it is generally not economical to remove the blood spot by hand, thus the blood spot is generally left in the commercial kirimi product. If the blood spot could be removed during processing, a higher quality kirimi product would result.\nAs can be seen from the above discussion, there exists a need for a method to remove the blood spot and any other internal organs remaining in the kirimi after processing. The present invention is directed toward fulfilling this need."}
-{"text": "Many therapeutic treatments are administered to a patient while they are sleeping or are attempting to fall asleep. While these treatments may achieve their intended result, they also often severely affect the quality of sleep that the patient gets while undergoing these treatments. These treatments often interrupt the patient's normal progression of sleep, causing transient arousals. While these arousals do not result in the awakening of the patient, they often pull patients from deeper stages or higher quality states of sleep. Patients often do not reenter these deeper stages of sleep for a relatively long period of time.\nIn some instances, a therapeutic treatment may cause numerous arousals. This fragments the patient's sleep and prevents the patient from reaching the deeper stages of sleep. Studies have shown that fragmented sleep results in excessive daytime sleepiness. This, in turn, is a direct contributor to many accidents, to a general feeling of lethargy, deterioration of cognitive performance, and/or daytime sleepiness, in the patient.\nOne example of therapeutic treatments causing sleep fragmentation is in the treatment of sleep disorders. Continuous Positive Air Pressure (CPAP) treatments are a primary remedy for a number of sleep disorders such as sleep apnea, hypopnea, and snoring. CPAP treatments consist of delivering a constant positive airway stream of air pressure into a patient's airway during sleep in order to keep the patient's airway from collapsing upon itself. State-of-the-art CPAP machines, often called auto-titration PAP (APAP) machines, automatically adjust the pressure of the delivered air in order to accommodate a patient's respiratory pattern. to the rapid changes of pressure in the patient's airway caused by the APAP machines. Another drawback of current state-of-the-art APAP machines is that they are subject to either false positives (such as when UAR and/or natural irregular breathing events are not pre-empted or do not occur, despite false detection of such and associated treatment control change) or false negatives (such as when genuine upper airway resistance (UAR) and/or related events are pre-empted or do occur but are not detected or responded to with treatment control change). This is due in part to the reliance of these machines on the correct interpretation of an inspiratory waveform and the inaccuracies related to the interpretation of the underlying waveform by the APAP machine. This can also be due to current state of the art gas delivery (or other treatment control such as pacemaker devices) devices inability to enable suitable algorithms to detect and adapt their computation detection sufficiently to pre-empt or predict the probability or onset likelihood of shallow breathing, UAR, arousals, and or associated sleep fragmentation or sleep quality deterioration.\nThe inspiratory waveform varies periodically for reasons not always associated with upper airway resistance. The use of inspiratory waveform as the primary or only means of detection of UAR-related events can cause remedial auto titration measures to be taken when none should be. This is particularly evident where the inspiratory waveform analysis technique does not employ an underlying time-course computational method. The time-course computational method refers to comparing a previous sequence of breaths (prestored from previous treatment session or stored from current session breathing data) or the current breath and comparing the variations or changes as an inferred measure of arousal or sleep fragmentation onset. Excessively rapid or excessively insensitive pressure changes often occur when an auto-CPAP machine tries to correct a normal non-UAR related event, or misses detecting the presence of subtle shallow breathing, hypopnea or UAR, respectively. It is believed that the primary cause of sleep fragmentation is the rapid pressure changes in the patient airway produced by the current APAP machines.\nIn addition to the above, studies have also suggested that some APAP machines are limited in their ability to accurately detect the onset or incidence of shallow breathing, mild hypopnea, or UAR events. This limitation is also possibly attributed to limitations of the machines in interpreting the wave form. Misdiagnosis of such mild hypopnea events results in increased UAR which in turn results in arousal and subsequent sleep fragmentation.\nCurrent state-of-the-art therapeutic devices do not optimally adapt to minimize arousals during therapy. Each patient's arousal threshold is affected by varying parameters, yet current state of the art devices do not have adaptive control algorithms that can adapt their treatment levels to accommodate a number of these varying parameters. These varying parameters include (but are not limited to) sleep history such as sleep deprivation or sleep propensity, physiological factors, psychological factors including (but not limited to) stress or anxiety, environmental factors including temperature; noise; lighting; vibration, factors such as varying threshold to arousals with changing age, drugs and alcohol effects to arousal thresholds and others.\nConsequently, in light of the inherent drawbacks in current therapeutic methods for administering treatments to patients who are sleeping or are attempting to sleep, there exists a need for an apparatus and method of monitoring for patient arousal and for adapting a therapeutic treatment to minimize arousal."}
-{"text": "The present invention relates to a process for the production of rumen bypass feed supplements. The process converts glyceride oils to their respective fatty acid calcium salts. In particular, the present invention relates to a process for the production of calcium salts of unsaturated fatty acids derived from fish oil. The calcium salts of the present invention, when fed to cattle, provide reproductive benefits, in particular, an increase in fertility as embodied in an increased rate of impregnation. The present invention therefore also relates to methods for providing such benefits in a ruminant.\nDairy cows must be impregnated once a year to maintain a lactation cycle in which milk is produced for ten months at a time with two month rest periods in between during which the cow is dry. Given the gestation period of a dairy cow, the objective is to impregnate the cow within 83 days after calving. The efficient management of a dairy herd thus requires that the cows be maintained at the peak of fertility to ensure re-impregnation within 83 days.\nAccordingly, there exists a need for nutritional supplements that promote dairy cow fertility. Fish oil fatty acids have become the focus of numerous research programs that seek to capitalize on their nutritional and physiological properties. WO 99/66877 discloses the use of omega-3 fatty acids of fish oil origin to increase fertility in animals including cattle. Among the omega-3 fatty acids disclosed are eicosapentaenoic acid (EPA) and docosahexaenoic acid (DHA).\nUnsaturated fatty acids, however, undergo hydrogenation to saturated fatty acids by microbial action in the rumen and must be fed to ruminants in a protected form. The most familiar form in which fatty acids in general are protected from microbial action in the rumen are the fatty acid calcium salts disclosed by U.S. Pat. Nos. 4,642,317; 4,826,694, 4,853,233; 4,853,233 and 4,909,138. This form of fatty acid protection is widely accepted in the dairy industry.\nFish oils have a glyceride content of 100%. That is, all of the fatty acids in fish oils are in the glyceride form. Fatty acid glycerides do not readily react to form calcium salts using the processes disclosed by the above-listed patents. For a product to be commercially feasible, glyceride levels below about 5 weight percent are desirable to produce a free-flowing and stable product.\nU.S. Pat. No. 5,382,678 discloses a process that reportedly can be used to prepare fatty acid calcium salts from feedstocks having glyceride contents as high as 40 weight percent, with the resulting product having a residual glyceride content of less than about 5 weight percent. Products with residual glyceride contents above 5 weight percent lack storage stability, and are susceptible to oxidation, post-heating, melting, subsequent product solidification, and a tendency to form lumps upon storage. Under industrial conditions, however, it has not been possible to consistently obtain residual glyceride levels below 5 weight percent once the initial glyceride content of the fatty acid feed stock is above about 25 weight percent when using the process of U.S. Pat. No. 5,382,678.\nHydrolyzing the glycerides to levels below 25 weight percent is not commercially feasible. Commercial omega-3 fatty acids in the free fatty acid form are so costly as to be commercially unfeasible. One can reduce the glyceride content of the fish oil starting material by blending it with a low glyceride content fatty acid feedstock, such as Palm Fatty Acid Distillate (PFAD), which has a glyceride content of about 15 to 20 weight percent.\nHowever, the quantity of PFAD that would have to be added to fish oil to reduce the glyceride content to levels commercially feasible for use with the process of U.S. Pat. No. 5,382,678 dilutes the concentration of desirable omega-3 fatty acid such as EPA and DHA to ineffective levels. That is, the levels of DHA and EPA in the resulting calcium salt are so low that quantities of calcium salt must be added to the daily feed ration at levels above what is considered acceptable by the dairy industry.\nTo be commercially viable, omega-3 fatty acid calcium salt feed supplements must have DHA and EPA concentrations high enough to confer the beneficial effects of these omega-3 fatty acids when quantities of the calcium salt are added to feed ration at levels considered acceptable to the cattle industry. Therefore, a need exists for a process by which calcium salts may be prepared from fish oils having high concentrations of omega-3 fatty acid with the calcium salts produced with reduced levels of unreacted glycerides in a free-flowing and stable form easily handled by customers.\nThis need is met by the present invention. It has now been discovered that fatty acid calcium salts having acceptable levels of residual glycerides can be prepared from high glyceride content starting materials by using elevated levels of calcium oxide, making it possible to prepare fatty acid calcium salts from feedstocks containing levels of fish oil effective to provide useful concentrations of omega-3 fatty acids in the finished product.\nTherefore, according to one aspect of the present invention, a method is provided for the preparation of fatty acid calcium salts, which includes the steps of:\n(a) providing a fatty acid feedstock having a glyceride content between about 30 and about 60% by weight;\n(b) adding to the feedstock from about 2 to about 3 equivalents of calcium oxide relative to the feedstock, so that a reactive admixture is formed; and\n(c) adding to the reactive admixture from about 2 to about 5 equivalents of water relative to the calcium oxide, so that the calcium oxide hydrates and neutralizes the fatty acids to form calcium salts.\nThe method of the present invention thus includes the use of feedstocks derived from fish oils diluted to glyceride contents between about 30 and about 60 weight percent with low glyceride content fatty acid feedstocks, such as PFAD. Other suitable sources of low glyceride content fatty acids include fatty acids from soy, cottonseed, corn and other vegetable fatty acid distillates, tallow, yellow grease or other animal or fish derived free fatty acid sources produced by deodorization, refining, hydrolyzation or other processes common in the fats and oil industry.\nThus, the method of the present invention obtains fatty acid calcium salts having useful concentration of omega-3 fatty acids and acceptable levels of residual glycerides that heretofore could not be obtained on a commercial scale using prior art manufacturing techniques. Therefore, according to another aspect of the present invention, fatty acid calcium salts are provided containing omega-3 fatty acids and residual glyceride levels below about 5 weight percent that are prepared by the method of the present invention. More specifically, a fatty acid calcium salt product is provided having a residual glyceride content below about 5 weight percent and containing from about 1 to about 10% by weight of EPA calcium salt and from about 1 to about 10% by weight of DHA calcium salt.\nThe DHA- and EPA-containing fatty acid calcium salts enhance the fertility of ruminants without using 100% glyceride content fish oil feedstocks. That is, beneficial results are obtained from feedstocks blended with fatty acids from sources other than fish oils.\nTherefore, according to still yet another aspect of the present invention, a method is provided for increasing fertility in a ruminant, in which the ruminant is fed an effective amount of the EPA- and DHA-containing fatty acid calcium salts of the present invention. The method of the present invention is particularly effective to enhance the fertility of female ruminants, especially dairy cows. Methods in accordance with the present invention begin feeding the supplements daily to a female ruminant from about 21 days before to about 28 days after parturition and feeding continue at least until conception occurs. The fertility enhancement obtained by the calcium salts of the present invention also includes a reduction in embryonic death in the months following conception. Therefore, methods in accordance with the present invention continue feeding the supplements to a female ruminant for at least 30 days, and preferably for at least 60 days after conception.\nThe above and other features and advantages of the present invention will become clear from the following description of the preferred embodiment.\nThe present invention provides a process by which high glyceride content fatty acid feedstocks may be converted to free-flowing powder or granular fatty acid calcium salt rumen bypass feed supplements, making it possible to prepare fatty acid calcium salts from fish oils, wherein the calcium salts contain useful and beneficial quantities of omega-3 fatty acids. The use of such high glyceride content fatty acid feedstocks represents a significant departure from conventional processes for the manufacture of fatty acid calcium salt feed supplements. The term xe2x80x9cglyceridexe2x80x9d as employed herein includes C10-C22 fatty acid monoglycerides, diglycerides, triglycerides, and any mixture thereof.\nIn a typical process according to the present invention, fatty acid feedstocks are added to a production vessel. The mixing should be accomplished in a kettle designed so that intensive and intimate contacting occurs between the calcium oxide and fat/oil admixture so that a homogeneous dispersion of the calcium oxide particles results. The types of internal mixing elements span a wide gap but would include those with propeller, turbine, plows with chopper blades, or preferably xe2x80x9cCowles-typexe2x80x9d dispersing blades as examples, but others may apply. These same devices would also be suitable for dispersing and homogenizing the water fraction into the fat/oil-calcium oxide admixture.\nFatty acid feedstocks are employed containing up to about 60 weight percent of the fatty acid content in the form of fatty acid glycerides. Glyceride levels between about 30 and about 60 weight percent are preferred.\nSuitable feedstocks include essentially any glyceride or glyceride derivative having a fatty acid profile determined to be nutritionally or physiologically beneficial to a ruminant. Beneficial fatty acid profiles are readily identified by those of ordinary skill in the art, and may be derived from any known source, inclusive of fatty acid sources of animal, vegetable or fish origin. This includes distillates and soap stocks of lard or tallow, vegetable oils such as canola oil, sunflower oil, safflower oil, rapeseed oil, soy bean oil, olive oil, corn oil, palm oil, and the like, and byproducts thereof, as well as fish oils and byproducts thereof.\nSuch fatty acid feedstocks typically contain from about 10 to about 100 weight percent of the fatty acid content in the form of fatty acid glycerides, from about 0 to about 90% by weight of free fatty aids, and less than 5% by weight of moisture, insolubles and unsaponifiables. When necessary, the glyceride content may be reduced to about 60% by weight and lower by adding fatty acid distillates such as PFAD to the feedstock or by pre-treatment to convert a portion of the glycerides to free fatty acids, either catalytically through the use of enzymes, including lipases, or by hydrolysis. Pre-treatment processes may also include processes that increase the level of desirable fatty acids, for example, cold acetone extraction may be used to increase the level of EPA and DHA in fish oil.\nThe present invention thus makes possible the preparation of fatty acid calcium salts from fish oils, which have a 100% glyceride content. According to one embodiment of the present invention, from about 15 to about 50 weight percent of fish oil is blended with from about 85 to about 50 weight percent of a fatty acid feedstock sufficiently low in glyceride content to provide a blend with a glyceride content of less than about 60% by weight. Blends of fish oils with PFAD within these weight ratios are included within the scope of the present invention. A blend containing from about 20 to about 35% by weight of fish oil is preferred.\nThe process of the present invention is particularly well suited for the preparation of fatty acid calcium salts containing beneficial levels of DHA and EPA from fatty acid feedstocks containing high levels of fish oil. Oils containing from about 7 to about 16% by weight DHA and from about 10 to about 17% by weight EPA are preferred.\nThe feedstock may also contain up to about 95% by weight of unsaturated C:16-C:22 fatty acids. Unsaturated fatty acid levels between 50 and about 80 weight percent are preferred. In general, unsaturated fatty acids having from 16 to 22 carbon atoms and from 1 to 6 double bonds are suitable for use with the present invention. Polyunsaturated fatty acids are preferred, with examples of desirable polyunsaturated fatty acids including fish oil-derived omega-3 and omega-6 fatty acids. Suitable fish oil sources include menhaden, herring, mackerel, caplin, tilapia, tuna, sardine, pacific saury, krill, and the like.\nIt may be necessary to heat the fatty acid feedstock to form a uniform, liquid admixture, depending upon the degree of saturation. A temperature up to about 175xc2x0 F. is suitable, with a temperature between about 120 and about 140xc2x0 F. being preferred.\nCalcium oxide is added to the fatty acid feedstock in the range from about 2 to about 3 equivalents relative to the fatty acid feedstock. A calcium oxide level between about 2.25 and 2.75 equivalents is preferred.\nWater is then added to hydrate the calcium oxide to its hydroxide form, creating a large amount of exothermic heat. The heat that is evolved usually is sufficient for the fatty acid neutralization reaction to proceed to completion, so that it is not necessary to supply heat to the reaction mixture. Between about 2 and about 5 equivalents of water relative to the calcium oxide is added to the reaction mixture, with between about 2.5 and about 3.5 equivalents being preferred.\nThe excess water is converted to steam by the exothermic heat generated and boils off rapidly. The reaction can be performed under atmospheric pressure, or under vacuum to draw off the steam.\nThe amount of time required for the reaction is between about 5 and about 60 minutes, and more typically between about 6 and about 15 minutes. The reaction is easily identified by the transformation of the admixture into a solid granular mass. Upon transfer from the reaction vessel, it can easily be processed into dry free-flowing particles.\nFish oil derived fatty acid calcium salt rumen bypass feed supplements of the present invention have a residual glyceride content below about 5 weight percent and contain from about 1 to about 10% by weight EPA calcium salt and from about 1 to about 10% by weight DHA calcium salt. Products containing from about 1.5 to about 9% by weight of EPA and DHA are preferred. The EPA- and DHA-containing feed supplements derived from fish oil/PFAD blends enhance the fertility of female ruminants. The present invention makes possible the commercially viable preparation of DHA- and EPA-rich calcium salt feed supplements for the enhancement of female ruminant fertility.\nThe present invention therefore includes fatty acid calcium salts having a fatty acid profile consistent with a profile resulting from blending from about 15 to about 50% by weight of fish oil with from about 85 to about 50% by weight of PFAD, wherein the fatty acid profile includes from about 1 to about 10% by weight of DHA and from about 1 to about 10% by weight of EPA.\nThe fatty acid calcium salt rumen inert feed supplements of the present invention may be conveniently fed to a ruminant admixed with a conventional ruminant feed. The feeds are typically vegetable materials edible by ruminants, such as legume hay, grass hay, corn silage, grass silage, legume silage, corn grain, oats, barley, distiller\"\"s grain, brewer\"\"s grain, soya bean meal and cottonseed meal. There is no particular lower limit of the calcium salt to be added to the ruminant feed, although in practice, amounts of the calcium salt below an amount that supplies 6 grams per day each of DHA and EPA are too small to provide significant fertility enhancement.\nThe fertility of female ruminants is enhanced when they are fed the EPA- and DHA-containing calcium salts starting as early as 21 days before parturition. While useful enhancement of fertility may be obtained by starting dietary supplementation at parturition or as late as 28 days following parturition, optimal results are obtained by earlier supplementation.\nThe dietary supplementation should continue daily at least until conception. However, because the DHA- and EPA-containing calcium salts of the present invention also enhance fertility by reducing embryonic death in the months following conception, the present invention also includes methods in which dietary supplementation is continued for at least 30 days and preferably for at least 60 days after conception. Beneficial results are obtained when dietary supplementation is continued up to 150 days after conception.\nThe DHA- and EPA-containing calcium salts of the present invention may be co-administered with additional quantities of other fatty acid calcium salts that are fed to ruminants for other purposes. The DHA- and EPA-containing calcium salts may be co-administered with a fatty acid calcium salt intended to supply energy to a high milk production ruminant, with a calcium salt having a fatty acid profile selected to modify the fatty acid profile of the milk fat or meat fat of the ruminant, or with both. Accordingly, the methods of the present invention for enhancing female ruminant fertility do not exclude the co-administration of other fatty acid calcium salt that do not contain DHA or EPA. One method in accordance with the present invention transitions the ruminant from the DHA- and EPA-containing calcium salts of the present invention to fatty acid calcium salts intended to supply energy to a high milk production ruminant once the fertility-enhancing benefits of the DHA- and EPA-containing calcium salts are no longer needed.\nThe following non-limiting examples set forth herein below illustrate certain aspects of the invention. All parts and percentages are by weight unless otherwise noted, and all temperatures are in degrees Fahrenheit."}
-{"text": "Individuals and enterprises are continually exposed to risk because of future events beyond their control. The outcome of those events can either positively or negatively impact on their wellbeing.\nIndividuals and enterprises should generally prefer not to face exposure to the possibility of adverse consequences, regardless of their perception of the likelihood of such events occurring. It is in their interest to consider foregoing `resources` they currently possess if doing so would reduce the possibility of being so greatly exposed to future outcomes.\nRisk can take many forms in view of the large range and type of future events which might result in adverse consequences. Risk can be categorised, in one instance, as `economic` in nature. Phenomena that constitute economic risk include: commodity prices, currency exchange rates, interest rates, property prices, share prices, inflation rates, company performance, and market event based indices.\nAnother characterisation of risk concerns `technical` phenomena. This can include things like the breakdown of an electricity generation plant, aircraft engine failure, and the damage to, or failure of, orbiting telecommunications satellites. The outcomes for each of these phenomena will be adverse for the users and/or supplier.\nOther forms of risk defy ready characterisation, such as weather-based (viz., rain damage or lightning strike), or other natural occurrences (viz., earthquakes or iceberg collision with sea-going vessels).\nThere are also less tangible risks associated with, for example, the emission of atmospheric pollutants or the disposal of intractable toxic wastes, in the sense that the future consequences are unknown, save that there is a notion, based on current information, that they could be adverse.\nThe capability to manage risk is more important today than it was in the past, and is likely to become ever-more important into the future, because there is an ever increasing exposure to a wider generic range of future phenomena beyond the control of individuals or enterprises. There is also a wider feasible range of possible future events, and greater uncertainty about the likelihood of occurrence, associated with any single future phenomenon viz., an increasing volatility.\nIt is also thought that individuals are now more risk-averse in recessionary times, when there are fewer available discretionary resources to trade-off to protect themselves from such adverse future events.\nIn the prior art, individuals and enterprises faced with `technical` risk have hedged against future outcomes by mechanisms such as the adoption of quality assurance practices, warranties, increased research and development activity (and associated intellectual property rights such as patents, utility models and registered designs), the purchase of modernised plant and equipment, and improved inventory, occupational health and safety and employer/employee relations practices.\nConsider a manufacturer of, say, integrated circuits (ICs), which has many clients wishing to purchase its ICs. The demand may result in a delay in delivery due to limited manufacturing capacity, thereby requiring advance production scheduling for orders already in-hand. Typically, the manufacturer will give a warranty to a purchaser as to measurable performance criteria for its ICs; if a batch does not perform to the specified criteria, the manufacturer is required by contract to replace that batch. That is, a purchaser may have no interest in obtaining monetary compensation for the poor quality ICs, as the purchaser needs the components for their own products. In that case, the `consideration` the warranty makes is the priority scheduling of a substitute batch of that type of IC, possibly displacing other scheduled production runs, or deferring delivery to another purchaser.\nSuch contractual arrangements are piece-meal in nature, and can only be struck between the manufacturer and each individual purchaser. They also leave the manufacturer exposed to claims from other customers whose orders are delayed by the re-scheduling. The manufacturer has no convenient mechanism available to it to hedge against such claims, perhaps by way of reserving production rights with another manufacturer, in lieu of unavailability of their own manufacturing facility.\nIn the face of such `economic` risk, it is known for individuals and enterprises to hedge against adverse outcomes by indirect means such as self-insurance, and directly by means such as futures contracts, forward contracts, and swaps.\nThere are disadvantages or limitations associated with such available economic risk management mechanisms. Particularly, they provide, at best, only indirect approaches to dealing with the risk management needs. The available mechanisms are relatively expensive, and provide limited phenomenon coverage, and therefore cannot meet the requirements of the party seeking to hedge against such wide-ranging future risk. The infrastructure and pay-out costs associated with switching between, say, a commodities market and a stock market are often prohibitive for entities small and large alike. As a consequence, entities find themselves saddled with obligations they have little control over and cannot escape.\nIn respect of the \"less tangible\" forms of risk, an example in the prior art of a form of management of that risk is that of `pollution rights` sold by the U.S. Environmental Protection Agency (EPA) in March 1993 for the atmospheric emission of sulphur dioxide. This was done by an auction of \"allowances\" permitting the release into the atmosphere. By the year 1995, any company or organisation emitting sulphur dioxide in the U.S. without enough allowances to cover their total emissions will face prosecution. This means polluters must either buy further allowances, or else modify or replace their plant and equipment to reduce these emissions. The EPA will regulate the total number of allowances able to be obtained. The existing allowances have already become a valuable tradeable `property` as between sulphur dioxide emitters, that is, even before the time when no further allowances will be able to be purchased.\nManagement techniques for the `less tangible` forms of risk are in their infancy. The existing forms indicate an emerging demand for systems and methods to enable effective management.\nSpecific examples in the prior art of patents relating to methods and apparatus which deal with various forms of risk management include British Patent No. 2 180 380, in the name of Merrill Lynch Pierce Fenner and Smith Incorporated, directed to an Automated Securities Trading Apparatus (corresponding to U.S. Pat. No. 4,674,004, and further related to U.S. Pat. Nos. 4,346,442 and 4,376,978). Other examples include U.S. Pat. No. 4,739,478 assigned to Lazard Freres and Co., directed to Methods and Apparatus for Restructuring Debt Obligations, U.S. Pat. No. 4,751,640 assigned to Citibank, N. A., directed to An Automated Investment System, and U.S. Pat. Nos. 4,752,877, 4,722,055, and 4,839,804 assigned to College Savings Bank directed to Methods and Apparatus for Funding Future Liability of Uncertain Cost.\nThe present invention comes about in view of the shortcomings of existing risk management mechanisms, and the perceived increasing importance of the management of risk relating to specified, yet unknown, future events.\nIn this sense, the invention is directed to something having economic value to individuals, enterprises and societies as a whole. Methods and apparatus that provide for the management of risk offer material advantages by, for example, minimising adverse future outcomes, providing both a form of compensation in the event of adverse future outcomes, and forms of risk management not otherwise supported or available in the prior art, and thus have value in the field of economic endeavour."}
-{"text": "This disclosure relates to powertrain systems and, more particularly, to cylinder deactivation powertrain systems.\nConventionally, applying cylinder deactivation on piston engine powertrains provides improved fuel economy over a combination of speeds and loads limited by excessive torque pulses. However, the limited speed-load range of those piston engine powertrains, where cylinder deactivation can be used, reduce the overall improvement in fuel economy to such a degree that few production vehicles employ cylinder deactivation. When cylinder deactivation is employed, one or more cylinders are fired and then skipped, or deactivated, and a time lapse occurs until the next cylinder fires. The resulting time lapse translates into roughness that can become a customer dissatisfier. This irregular or staggered torque output reduces the fuel economy efficiency otherwise available from cylinder deactivation. In addition, the exhaust emissions generally increase due to the inconsistent fuel delivery and combustion.\nAccordingly, there exists a need to selectively fire cylinders without causing objectionable torque pulses and increasing emissions.\nThe drawbacks and disadvantages of the prior art are overcome by the embodiments of the powertrain system described herein. The powertrain system comprises an engine coupled to a drive shaft, and a torque-smoothing device that is configured to sometimes assist the drive shaft by adding energy, and at other times to move energy from the drive shaft to the energy storage device. The engine is also equipped with a cylinder deactivation system configured to deactivate one or more engine cylinders, and a direct cylinder injection system configured to deliver fuel to firing cylinders. An engine management system is configured to provide control, diagnostic and maintenance operations for the powertrain system.\nThe method for operating the internal combustion engine comprises activating the engine, and deactivating one or more engine cylinders. When the engine generates an inconsistent torque output, a torque-smoothing system is activated to smooth the torque output of the engine."}
-{"text": "The present invention relates to an automatic focus control device having a plurality of focus detecting regions.\nConventionally, an automatic focus detecting device having a single focus detecting region, in which a defocus amount is corrected in accordance with assembly errors of an optical system, has been proposed in, for example, Japanese Patent Laid-Open Publication No. 126517/1984. However, an automatic focus detecting device having a plurality of focus detecting regions, in which a defocus amount is corrected for each of the focus detecting regions, is not known.\nAs shown in FIG. 7, in the case where automatic focus detection is performed in three focus detecting regions, i.e. a horizontally elongated focus detecting area A disposed at a center of a photographing field S and a pair of vertically elongated focus detecting regions B and C disposed at opposite sides of the photographing field S, CCD light-receiving element arrays e1, e2 and e3 corresponding to the focus detection regions A, B and C, respectively are provided on one chip as shown in FIG. 4. If only one CCD light-receiving element array, for example, only the CCD light-receiving element array corresponding to the central focus detecting region A is provided, positional adjustment of the CCD chip in the direction of its optical axis can be performed by performing positional adjustment of only the central portion of the CCD chip in the direction of the optical axis. However, in the case where a plurality of, for example, three CCD light-receiving element arrays are provided as described above, parallelism of the CCD chip plays an important role. Namely, when parallelism of the CCD chip is poor, deviation in the direction of the optical axis occurs at opposite sides of the CCD chip. When parallelism of the CCD chip is poor, deviation in the direction of the optical axis occurs at opposite sides of the CCD chip. Even in the case of only one CCD light-receiving element array, the deviation also takes place at opposite sides of the CCD chip but assumes a small value. Meanwhile, also in the case where a plurality of the CCD light-receiving element arrays corresponding to the focus detecting regions, respectively are not provided on a single chip, deviations of the assembled CCD light-receiving element arrays in the direction of the optical axis are different from each other, so that it becomes necessary to correct the respective focus detecting regions in the direction of the optical axis.\nAs shown in FIG. 6a, a focus detecting device has been widely used in which an image formed by a light-receiving lens a is re-formed by a pair of reimaging lenses d1 into first and second images on a light.TM.receiving element array el arranged in a straight line and a distance between the first and second images is detected by the light-receiving element array el such that a focusing state is detected. In this known focus detecting device, when the, distance between the first and second images assumes a predetermined value, a decision of an in-focus state in which the image is formed on a predetermined focal plane is made. Meanwhile, when the distance between the first and second images is smaller and larger than the predetermined value, decisions of front and rear focus states in which the image is formed forwardly and rearwardly of the predetermined focal plane are made, respectively such that deviation amounts of the distance between the images from that of the in-focus position are outputted as defocus amounts. At the time of focus detection, the distance j between the first and second images is obtained and the defocus amount .DELTA..epsilon. is calculated by multiplying the distance j by a predetermined coefficient s. The above described Japanese Patent Laid-Open Publication proposes correction of the coefficient s since the coefficient s for the front focus state is different from that for the rear focus state.\nIn the CCD chip of FIG. 4, it is considered that if a length of the CCD light-receiving element arrays e1 and e3 corresponding to the focus detecting regions B and C disposed at opposite sides of the photographing field, respectively is made smaller than that of the CCD light-receiving element array e2 corresponding to the focus detecting region A disposed at the center of the photographing field, a focus detecting module can be made smaller in size by a difference between the length of the CCD light-receiving element arrays e1 and e3 and that of the CCD light-receiving element array e2 and a time period required for effecting data dump from the CCD light-receiving element arrays el and e3 can be reduced. To this end, as shown in FIG. 4, a pair of reimaging lenses d1 and a pair of reimaging lenses d3, which are, respectively, disposed at opposite sides of the CCD chip, have distances Dd1 and Dd3 smaller than a distance Dd2 between a pair of reimaging lenses d2 disposed at a center of the CCD chip, so that the numbers of the elements required for the CCD light-receiving element arrays e1 and e3 are determined accordingly. However, in this case, such a problem arises due to the difference between the distances Dd1 and Dd3 and the distance Dd2 that a coefficient for converting a distance between the first and second images into a defocus amount on the CCD light-receiving element array e2 is different from a coefficient for converting a distance between the first and second images into a defocus amount on the CCD light-receiving element array e1 and a coefficient for converting a distance between the first and second images into a defocus amount on the light-receiving element array e3.\nIn FIG. 6a, the distance De1 between the first and second images on the light-receiving element array e1 changes according to changes of a distance Dd1 of the reimaging lenses d1. Therefore, in the case where the reimaging lenses d1 are made of, for example, plastic, the distance Dd1 of the reimaging lenses d1 varies in response to temperature change, so that the distance De1 between the first and second images on the light-receiving element array e1 changes, thereby resulting in also change of the defocus amount. Thus, Japanese Patent Laid-Open Publication No. 235110/1985 discloses that temperature of the focus detecting device is detected such that the defocus amount is corrected in accordance with the detected temperature of the focus detecting device.\nMeanwhile, if temperature rise occurs in the case where as shown in FIG. 4, a pair of the reimaging lenses d1 and a pair of the reimaging lenses d3, which are, respectively, disposed at opposite sides of the CCD chip, have the distances Dd1 and Dd3 smaller than a distance Dd2 of a pair of the reimaging lenses d2 disposed at a center of the CCD chip, so that the numbers of the elements required for the CCD light-receiving element arrays e1 and e3 are determined accordingly, amounts of change of the distances between the two images on the CCD light-receiving element arrays e1, e2 and e3 upon temperature rise are different from each other due to the differences among the distances Dd1, Dd2 and Dd3 even if the reimaging lenses d1, d2 and d3 are made of an identical material. Meanwhile, in the case where the reimaging lenses are not made of an identical material, the reimaging lenses have different coefficients of thermal expansion, so that amounts of change of the distances between the images on the CCD light-receiving element arrays e1, e2 and e3 upon temperature rise are different from each other for the respective focus detecting regions."}
-{"text": "1. Field of the Invention\nThe present invention relates generally to a surgical instrument for insertion into human tissue and a method for use thereof.\n2. Description of the Prior Art\nInstruments and methods exist to penetrate and dilate tissue. One problem with conventional instruments is that when the instrument includes multiple components, it is often difficult to efficiently and reliably secure the components to one another. Another problem is that when multiple components have multiple handles, the disengagement of one component from another is more cumbersome, resulting in surgical inefficiency.\nIn surgical procedures involving sequential dilation of tissue, it is typical to use multiple instruments in succession to dilate the tissue to meet a desired surgical objective. For example, a guide wire is often inserted first, followed by a larger diameter stylet, which is then removed, followed by yet another instrument having a larger diameter to increase the diameter of the opening. Successive insertion and removal of multiple instruments is a time-consuming process that can hamper the surgical procedure. Additionally, the positioning of separate smaller instruments within larger instruments may lead to imprecise results as it is more difficult to maintain the instruments coaxial with one another, or limit the depth of insertion of the instruments relative to one another.\nAccordingly, there exists a need for a surgical instrument having a configuration so that multiple components may be individually or in combination engaged to a handle for easier insertion and removal of the surgical instrument and its components from a patient. There also exists a need for a surgical instrument capable of performing sequential dilation without the need to successively insert and remove a multitude of individual instruments."}
-{"text": "1. Field of the Invention\nThe present invention relates to a method and to a circuit arrangement for determining the quality of virtual connections extending via input lines and output lines of an asynchronous transfer mode (ATM) switching equipment, whereby respective message cells are transmitted according to an asynchronous transfer mode process during the course of these virtual connections and each of these message cells thereby has a cell header available for identifying the respective virtual connection for the input line or, respectively, output line coming into consideration.\n2. Description of the Prior Art\nThe problem generally arises in telecommunication switching systems of monitoring the quality of defined connections in order to be able to provide subscribers with appropriate information, for example, given complaints. The operators of public or private switching equipment, however, also have the need for or, respectively, raise the requirement of being able to measure and monitor the quality of the switching equipment. Quality measurements of this type also enter into the maintenance concepts in order to be able to, for example, avoid predictable deteriorations in quality beyond an allowable tolerance limit.\nGiven a method and a circuit arrangement of the type initially set forth, it is therefore an object of the present invention to provide a way to determine the quality of virtual connections on the basis of the message cells transmitted during the course of such virtual connections.\nIn a method of the type set forth above, this object is achieved, according to the present invention, in an improved method for determining the quality of virtual connections extending via input lines and output lines of an ATM switching equipment in which respective message cells are transmitted according to an asynchronous transfer mode during the course of these virtual connections and each of the message cells thereby has a cell header available with which the respective virtual connection for the input line or, respectively, output line coming into consideration is identified, and is particularly characterized in that the plurality of message cells accepted and output are separately acquired in a call-associated manner in the ATM switching equipment on the basis of the cell headers of the message cells accepted via the input lines or, respectively, output to the output lines, and in that, in response to the cleardown of the respective virtual connection, the loss rate of the message cells is calculated as a measure for the quality of the respective virtual connection, being determined from the previously, separately-acquired plurality of accepted and output message cells.\nThe present invention yields the advantage that the quality of the individual virtual connections can be continuously determined on the basis of the transmitted message cells without significantly dynamically loading an ATM switching equipment for this determination.\nAccording to a particular feature of the invention, the aforementioned improved method is further particularly characterized in that a mean loss rate is formed as quality measure for the ATM switching equipment on the formation being from the identified loss rates of the virtual connections established within a defined time interval. Also, the improved method is particularly characterized in that the loss rate identified for the respective virtual connection is involved in the calculation of charges. This feature and development yields the advantage of a simple determination of the quality of an ATM switching equipment.\nThe advantage of the improved circuit arrangement set forth above is that only the message quantity forwarded by an ATM switching equipment to a subscriber equipment, in fact, enters into the calculation of charges. Message cells repeatedly transmitted within an ATM switching equipment or, respectively, misrouted message cells due, for example, to transmission malfunctions are left out of consideration for calculating charges.\nThe circuit arrangement for the implementation of the method set forth above is characterized in that a plurality of counters corresponding in number to the plurality of virtual connections extending via the input lines and the output lines is assigned to the input lines and to the output lines, in that the plurality of counters can be placed into a defined initial counter reading proceeding from a central control device during the respective course of the setup of the allocated virtual connection, in that the individual counters are individually drivable given the appearance of message cells on the allocated input line or, respectively, output line, being individually drivable according to the prescription of the cell header respectively contained in the message cells, and in that the momentary counter reading of a counter is variable by a counting unit proceeding from a defined initial counter reading, the reading being variable with each drive, and in that the momentary counter readings of the counter assigned to a virtual connection can be interrogated at least at the cleardown of the connection by a central control device for calculating the loss rate.\nThe advantage of the circuit arrangement is thereby that the registration of the message cells transmitted during the course of virtual connections occurs in peripheral equipment so that the central control device of an ATM switching equipment is largely relieved of the determination of the quality of the virtual connections.\nAccording to other features and developments of the invention, the circuit arrangement is particularly characterized in that the momentary counter readings of the counter can be additionally interrogated at fixed time intervals. Also, the arrangement is particularly characterized in that the counters are fashioned as periodically-circulating counters and in that a message signal is separately offered to the central control device before a cleardown of the respective virtual connection each time a defined counter reading is reached. The advantage of these features and developments is that the counters provided for the registration of message cells require only a relatively low counting capacity and, therefore, the required circuit expense can be kept low."}
-{"text": "1. Field of the Invention\nThe present invention relates to reflective warning devices which are useful for alerting motorists to the presence of other vehicles and road hazards.\n2. Description of the Prior Art\nThe prior art has previously provided warning devices of both an active and passive nature which act to warn of the presence of a dangerous condition, particularly a road hazard associated with moving or stationary vehicles. Active visual devices particularly include lights disposed on a vehicle or on a stationary object intended to warn of a dangerous condition, it having been previously appreciated that a flashing light is more effective in alerting observers to the presence of potential dangers. However, active visual devices must be operated by a source of power such as a battery or the like which causes the light to repetitively flash on and off. As is obvious, such active devices require a substantial amount of energy for operation and are also costly. Passive visual warning devices, particularly reflective devices wherein light from vehicle headlights or other sources are reflected to observers, are also known in the art. Such reflective devices have even been caused to provide a flashing effect by intermittently obscuring at least portions of a reflective surface. For example, LaLonde, in U.S. Pat. No. 3,528,721, causes a flashing effect of a reflective device by mounting the device against the wheel of a bicycle so that rotation of the wheel causes the reflective device to oscillate. Priest, in U.S. Pat. No. 3,551,024, discloses a flashing retroreflective reflector in which a convex lens focuses incoming light on a retroreflective medium, the device being oscillated such that retroreflection is obtained only when incoming light is incident along the axis of the lens. Hammer, in U.S. Pat. No. 2,869,424, provides a passive warning signal device wherein portions of a reflective surface are intermittently obscured by spring mounted masking plates, the masking plates being oscillative due to motion of the vehicle carrying the reflective device. Klaenhammer, et al., in U.S. Pat. No. 4,023,888, provides a reflective structure having portions thereof intermittently obscured by a swinging masking member, the swinging member being driven by a motor. However, the prior art has not provided a passive visual warning reflector device capable of producing a \"flashing\" visual signal on mounting of the device to either a stationary or potentially movable object. A free-swinging portion of the present device has means formed thereon which are actuable by air currents to be operable both when the structure on which the device is mounted is moving and also when the structure is stationary, as long as air is circulating about the device when the structure is stationary."}
-{"text": "Field\nThe technology described herein relates generally to a liquid crystal display (LCD), and more particularly, to an LCD including a plurality of pixels displaying a color image.\nDiscussion of the Background\nA typical liquid crystal display (LCD) includes a pixel electrode, a common electrode, a liquid crystal (LC) layer including liquid crystal molecules configured to be tilted by an electric field that is generated by the pixel and common electrodes, and a backlight unit configured to emit light to the LC layer.\nSuch an LCD device displays a desired image by controlling intensity of the electric field generated by the pixel and common electrodes and then adjusting the transmittance of light that is emitted from the backlight unit to be transmitted through the LC layer.\nThe LC layer included in the LCD may have a phase delay characteristic that is intrinsic to the liquid crystals themselves, and accordingly, in the LCD, maximum transmittance of light transmitted through the LC layer has a characteristic according to each wavelength.\nThe above information disclosed in this Background section is only to enhancement of understanding of the background of the inventive concept, and, therefore, it may contain information that does not form the prior art that is already known to a person of ordinary skill in the art."}
-{"text": "The present disclosure relates to a sheet conveying device for conveying a sheet member, and an image reading device and an image forming apparatus including the sheet conveying device including a conveying roller configured to be rotatable in a forward rotational direction and a reverse rotational direction.\nThe image reading device such as a scanner is provided with an automatic document feeder (hereinafter referred to as \u201cADF\u201d) for automatically conveying a document sheet set in a sheet feeding tray. The ADF includes a feeding roller for feeding inside the document sheet set in the sheet feeding tray, and a discharging roller for discharging the document sheet to an external discharge tray. The document sheet conveyed to a reading position by the feeding roller is discharged to the discharge tray by the discharging roller.\nIn an image reading device including the ADF as described above, conventionally, there is known a stopper member for preventing the document sheet in the discharge tray from being drawn into the discharging roller, when the discharging roller is rotated in a rotational direction opposite to that at a time of discharge. The stopper member is configured to prevent a rear end of the document sheet from coming in contact with a surface of the discharging roller for discharging the document sheet."}
-{"text": "This invention relates to an adapter for a discharge lamp. The adapter is particularly suited for use with a discharge lamp in which a discharge tube is disposed substantially in a plane and a lamp support housing is substantially surrounded by the discharge tube.\nU.S. Pat. No. 4,241,386 discloses an adapter for a discharge lamp having a circular discharge tube, which surrounds a circular zone. The discharge lamp has a plug portion disposed on the peripheral region of the circular zone. Accordingly, this known adapter also comprises a socket portion matching the plug portion of the circular discharge lamp. The socket portion must be positioned in a peripheral region of the adapter in order to accommodate the circular discharge lamp.\nMost mounting sockets receiving such circular low pressure discharge lamps (also called fluorescent lamps) have a socket portion which is situated in a peripheral area of the mounting socket, corresponding to the peripheral position of the plug portion of the circular discharge lamp.\nU.S. Pat. Nos. 4,549,251 and 4,458,301 disclose a discharge lamp with a discharge tube where the discharge tube is bent in a shape so as to surround a substantially planar zone. The discharge lamp includes a lamp support housing. The lamp support housing is positioned in an substantially central area of the zone defined by the discharge tube. The lamp support housing comprises a plug portion with contact pins providing the electrical connection of the discharge lamp.\nIt would be desirable if discharge lamps of the type described in U.S. Pat. Nos. 4,549,251 and 4,458,301 were connectable to mounting sockets originally designed for circular lamps, i. e. for mounting sockets which have their socket portion in a peripheral area.\nU.S. Pat. No. 5,471,375 teaches a mounting socket for use with a discharge lamp having a similar layout as those described in U.S. Pat. Nos. 4,549,251 and 4,458,301. The mounting socket comprises a socket portion matching the electric contacts of the discharge lamp, and also comprises a plug portion matching a known screw-type socket for incandescent lamps. Both the socket portion and the plug portion of the mounting socket are positioned in a central area of the discharge lamp when the discharge lamp is inserted in the mounting socket. However, when the mounting socket is placed in the screw-type socket of the incandescent lamp, and the discharge lamp is inserted in the mounting socket, the plane of the discharge lamp socket is above the level of the screw-type plug portion of the mounting socket. For this reason, a discharge lamp inserted in the mounting socket may interfere with the proper positioning of a lamp cover originally designed to cover a lamp inserted in the lower socket. Thus there is a particular need for a structure in which the level of the discharge lamp plug portion is not much higher than the plug portion of the mounting socket, and the plane of the discharge lamp is substantially the same as the plane of the original socket.\nIn an exemplary embodiment of the invention, an adapter for a discharge lamp comprises a socket portion suitable for receiving electrical contacts of a first discharge lamp. The first discharge lamp includes a discharge tube, which is disposed substantially in a planar zone having a principal plane. The socket portion has an associated socket contacting plane. The associated socket contacting plane is defined as the principal plane of the discharge tube of the first discharge lamp, when the latter is inserted into the socket portion of the adapter.\nThe adapter also comprises a plug portion which is insertable into an external female socket. The external female socket is of the type which is suitable for receiving a second discharge lamp, where the second discharge lamp has a discharge tube with a planar tube configuration. The plug portion of the adapter has an associated plug contacting plane. This plug contacting plane is defined by the principal plane of the tube configuration of the second discharge lamp, when the latter is inserted into the external female socket. The adapter further comprises a connecting portion for connecting the socket portion and the plug portion. The connecting portion provides a substantially rigid connection between the socket portion of the adapter and the plug portion of the adapter so that the socket contacting plane substantially coincides with the plug contacting plane.\nThe adapter is particularly suitable for discharge lamps which have a discharge tube surrounding a lamp support housing so that the lamp support housing is positioned in an substantially central area of the plane defined by the discharge tube. Such a lamp will be termed, folded lamp,, hereinafter because the ends of the discharge tube are folded inwards into the lamp support housing. The connecting portion of the adapter may be sized so that the plug portion is positioned in an substantially peripheral area of the plane defined by the discharge tube when the plug portion of the discharge lamp is inserted in the socket portion of the adapter. This allows discharge lamps with central plug portions to be used with mounting sockets and lamp housings which were designed for circular discharge lamps. Due to the fact that the plug contacting plane of the adapter coincides with the socket contacting plane, the principal plane of the discharge tube of the folded lamp will be practically in the same plane as the plane of the circular lamp. This means that the lamp covers originally designed for circular discharge lamps will be readily applicable with folded lamps as well."}
-{"text": "1. Field of the Invention\nThe present general inventive concept relates to a carrier recovery apparatus usable with a vestigial side band (VSB) receiver and a method thereof. More specifically, the present general inventive concept relates to a carrier recovery apparatus using not only a pilot signal, but also upper and lower sidebands of a received VSB signal, and a method thereof.\n2. Description of the Related Art\nIn order for a data receiver to accurately demodulate VSB-modulated data, it is necessary to minimize frequency offset and high levels of phase noise (jitter) generated from a tuner or an RF (radio frequency) oscillator used for data receiving. This procedure is called \u2018carrier recovery\u2019.\nA digital broadcast system based on a VSB modulation method in conformity with the standards of the Advanced Television Systems Committee (ATSC) uses a pilot signal in a transmitting signal for carrier synchronization. Here, the pilot signal is a signal loaded on a carrier during the transmission for accurate carrier recovery.\nFIG. 1 is a block diagram illustrating a conventional carrier recovery circuit 100 for phase detection. Referring to FIG. 1, the conventional carrier recovery circuit 100 comprises a multiplier 101, a pilot detector 103, a phase detector 105, a loop filter 107, and a numerically controlled oscillator (NCO) 109.\nA received VSB signal is digitalized by an analog-to-digital converter (ADC) (not shown), and output as a baseband signal by the multiplier 101.\nThe pilot detector 103 detects a pilot signal from the baseband signal, and the phase detector 105 reads phase information of the pilot signal. There are many methods for reading phase information of the pilot signal, and a suitable method is selected depending on the application. The phase information read by the phase detector 105 goes through the loop filter 107, is converted to a frequency component through the NCO 109, and is multiplied with the received VSB signal by the multiplier 101 to output the baseband signal.\nThe above-described procedure is repeatedly performed according to a feedback operation until a phase error of the pilot signal becomes zero.\nHowever, if the pilot signal in the wireless environment is corrupted, a received signal cannot be recovered properly. The unwanted pilot signal corruption occurs often in a multi-path environment, and eventually causes performance degradation of the VSB receiver. This problem can occur in any American digital broadcasting system using the VSB method."}
-{"text": "U.s. pat. No. 2,953,460 PA1 U.s. pat. No. 3,708,255 PA1 U.s. pat. No. 3,882,768"}
-{"text": "The present invention relates to a mobile terminal that can integrally move between an infrastructure network and an ad hoc network and to a method of controlling the same.\nAn infrastructure network in which networks such as internets are interconnected via a relay node such as a router as well as an ad hoc network which has no relay node and is a network temporarily configured of only terminals are well known as aspects of a network.\nIn order to establish communications by connecting a mobile terminal to a network, the mobile terminal must have (1) a terminal address used in a network to be connected, or a network address of the network itself and (2) an opposite communication party's address.\nWith respect to the item (1), the dynamic host configuration protocol, DHCP, (for example, refer to Douglas E. Comer, \"Internetworking with TCP/IP Volume I principles, protocols and architecture\", Third edition, 1995, Prentice-Hall, Inc.) is well known with an internet being an example of an infrastructure network. Moreover, with respect to the item (2), the domain name system, DNS, (for example, refer to Douglas E. Comer, \"Internetworking with TCP/IP Volume I principles, protocols and architecture\", Third edition, 1995, Prentice-Hall, Inc.) is well known. In the DHCP, a DHCP server is set up that holds a terminal address or infrastructure network address to be allocated to a terminal. Necessary values such as terminal or network addresses for connection are requested to the DHCP server when a terminal is connected to a network. The DHCP server provides parameters to the terminal without troubles due to duplicate terminal addresses, based on the request from the terminal. FIG. 15 depicts such an operation. When the terminal 1 and the terminal 2 are connected to a network, they respectively broadcast DHCP Requests containing desired information. FIG. 15 shows an example in which the terminals 1 and 2 request a terminal address. In response to the DHCP Request message, the DHCP server transmits a DHCP ACK message containing supply information to the request sources. In FIG. 15, the DHCP server transmits two DHCP ACK messages, one containing an address 1 and the other containing an address 2, to avoid overlapping of a terminal address at the terminals 1 and 2.\nIn the internet, the DNS (Domian Name System) is used to detect the correspondence between the terminal name of a terminal and the terminal address. The DNS introduces a DNS server that manages the correspondences between the terminal's names of terminals to be managed and the terminal addresses and predetermines the procedure of making inquiries from a terminal to the DNS server and the procedure of making inquiries between servers. In most cases, the address of the DNS server managing a terminal is set to the terminal itself. FIG. 16 depicts the case where the terminal 1 detects the terminal address of the terminal 2. The terminal 1 transmits a DNS Query message requesting the terminal address of the terminal 2 to the DNS server (the DNS server 1 in FIG. 16) registered. The DNS server 1 inquires the DNS server (the DNS server 2 in FIG. 16) managing the terminal address of the terminal 2. Then the DNS server 1 notifies the terminal 1 of the address 2 (DNS Reply message) when obtaining the terminal address of the terminal 2 (address 2).\nThe ad hoc network is configured of only terminals and does not have a server such as DHCP server or DNS server used in the infrastructure network. For this reason, even if the same procedure as that in the infrastructure network are used, communications cannot be accomplished by connecting the terminal to an ad hoc network. In order to accomplish communications according to the same procedure as that in the infrastructure network, a server function may be added to one of terminals connected to an ad hoc network. However, since this approach requires a special terminal with a server function, it is considerably poor in versatility. In order to improve the versatility, a procedure different from that in the infrastructure network may be defined in the ad hoc network so that the procedure can be selectively used according to an ad hoc network or infrastructure network to be connected. However, in the case of this procedure, the terminal must be reset when it moves between an ad hoc network and an infrastructure network according to the procedure for an infrastructure network and the procedure for an ad hoc network to be introduced thereto."}
-{"text": "A wide variety of fibrous materials have been employed in tobacco smoke filter elements. However, the choice of materials for use in production of such filters has been limited because of the need to balance various commercial requirements. A very important property of a tobacco smoke filter is obviously its filtration efficiency, i.e., its ability to remove selected constituents from the tobacco smoke. However, the range of filtration efficiency has had to be compromised in order to satisfy other commercially important factors such as resistance to draw, hardness, impact on taste, and manufacturing costs.\nCellulose acetate has long been considered the material of choice in the production of tobacco smoke filters, primarily because of its ability to provide commercially acceptable filtration efficiency, on the order of about 50%, without significantly detracting from the tobacco taste, low resistance to draw, and filter hardness desired by the majority of smokers.\nA significant component of the commercially desirable \"taste\" is provided by the standard plasticizers utilized in the production of filter elements from cellulose acetate fibers, usually triethylene glycol acetate or glycerol triacetate (\"triacetin\"). In conventional cigarette filter manufacturing, the plasticizer is commonly applied to the cellulose acetate fiber by spraying or wicking using art-recognized techniques. The tendency of the plasticizer to migrate toward the center of conventional cellulose acetate fibers reduces the level of plasticizer at the fiber surface, minimizing its taste-enhancing capability and limiting the shelf life of plasticized tow fibers before being processed into filter rods. The plasticizer is therefore usually added to the tow during the manufacture of the filter rods.\nCellulose acetate fiber plasticized in this manner and wrapped with paper into rod-like forms become bondable at the fiber contact points, enabling the formation of relative self-sustaining, elongated filter rods in two to four hours. This process can be accelerated by the application of gases at elevated temperatures simultaneously with the formation of the filter rod. Filter rods produced in this manner provide a tortuous path for the passage of tobacco smoke when discrete lengths of such material are utilized as tobacco smoke filter elements.\nFiltration efficiency can be increased significantly through the use of small fibers which provide increased fiber surface area at the same weight of fiber. Solvent spun cellulose acetate fiber is commercially available only in fiber sizes down to 13 microns in diameter. To obtain finer cellulose acetate fiber, e.g., 10 microns or less, melt spinning of plasticized cellulose acetate resin would be required; however, the level of plasticizer necessary to directly spin such fine cellulose acetate fibers would render the resultant fibers very weak and commercially useless. Melt spun cellulose acetate of a larger diameter, which would require less plasticizer, would have to be drawn and crimped to produce such fine fibers for use in tobacco smoke filters. Unfortunately, melt spun cellulose acetate fibers can only be commercially drawn at relatively low draw ratios before the fibers break during processing. The inability to form and process very fine fibers of cellulose acetate places practical limits on the filtration efficiency capabilities of this material in the production of tobacco smoke filters.\nFurther, and very important commercially, by comparison with other polymeric materials such as the polyolefins, cellulose acetate is relatively expensive, costing, for example, on the order of more than three times as much as commercially available polypropylene in resin form. While attempts have been made to utilize other less expensive and more easily processed polymeric materials such as polypropylene in lieu of cellulose acetate in the manufacture of tobacco smoke filters, such efforts have been almost universally abandoned on a commercial level, primarily because of the undesirable impact of such materials on the taste properties of tobacco smoke. Also, such use is generally limited by the inability to easily bond the fibers in order to obtain the desired filter hardness at required resistance to draw.\nAnother problem with commercially available tobacco smoke filters, particularly cigarette filters, currently on the market is the difficulty in disposing of such materials after use. By bonding highly crimped cellulose acetate fibers at their contact points, conventional cigarette filters are designed to provide a significant volume of interstitial space for the passage of smoke. The bonded contact points of such filter elements degrade very slowly under normal environmental conditions resulting in high volume, long life, environmentally undesirable litter."}
-{"text": "1. Field of the Invention\nThe present invention relates to a demodulation circuit and a demodulating method for demodulating a digital-modulated signal.\n2. Description of the Related Art\nFIG. 1 is a block diagram of a common QPSK demodulation circuit. A QPSK-modulated signal is received by a tuner not shown in the attached drawings, and converted into an intermediate frequency signal IFin. The intermediate frequency signal IFin is converted by an A/D converter 12 into digital data at the timing synchronous with the sampling clock with a predetermined frequency, and further converted by an IQ separator 13 into a baseband signal formed by an I signal and a Q signal.\nThe IQ separator 13 includes multipliers 13a and 13b for multiplying the output of the A/D converter 12 by a sine wave and a cosine wave having the frequencies equal to the center frequency of the intermediate frequency signal IFin. Low pass filters (LPF) 14a and 14b remove the upper frequency components of the output of the multipliers 13a and 13b, and the output is transmitted to interpolaters 15a and 15b. \nEach of the interpolaters 15a and 15b is constituted by an FIR (finite impulse response) filter, a thinning circuit, etc., and the data of the original sampling point (symbol timing) is obtained by interpolation from the received data sampled by the A/D converter 12.\nEach of the Root Nyquist filters 16a and 16b is constituted by a low pass filter, and restricts the band of the output signal of the interpolaters 15a and 15b. \nA rotor 17 is constituted by a butterfly circuit, etc., advances or delays the phase of a regenerated carrier wave according to the phase difference signal output from a carrier recovery circuit 18, and thereby allows the regenerated carrier wave to be synchronous with the carrier wave of a received signal.\nThe carrier recovery circuit 18 calculates the phase difference between the signal point on the I and Q phase plane obtained from the I signal and the Q signal and the normal signal point, and outputs the signal in the phase difference decreasing direction to the rotor 17.\nA timing recovery circuit 19 calculates the amount of shift of the sampling timing, and output an interpolating tap coefficient to the interpolaters 15a and 15b. \nFIG. 2 shows the configuration of a timing recovery loop. The timing recovery loop is constituted by the interpolaters 15a and 15b and the timing recovery circuit 19.\nThe timing recovery circuit 19 includes a phase comparator 21, a loop filter 22, a numerical controlled oscillator 23, a thinning control unit 24, and a tap coefficient arithmetic unit 25. The interpolaters 15a and 15b includes Fir filters 26a and 26b, and thinning circuits, 27a and 27b. \nThe phase comparator 21 determines whether the value of a received signal at each sampling timing indicates a monotonous increase, a monotonous decrease, or none of them, and outputs a signal depending on a discrimination result.\nThe loop filter 22 is a completely secondary loop filter, and includes a multiplier 28 constituting a low pass filter, a multiplier 29 constituting a high pass filter, an adder 30 for adding the output of the multiplier 28 to the output of the integrator 31, a limiter 32, an integrator 31, and an adder 33 for adding the output of the integrator 31 to the output of the multiplier 29. The integrator 31 is constituted by a flip-flop, etc. The coefficient \u03b1 provided for the multiplier 28 is a gain adjustment coefficient for a low pass filter, and the coefficient \u03b2 provided for the multiplier 29 is a gain adjustment coefficient for a high pass filter.\nThe numerical controlled oscillator (NCO) 23 includes a flip-flop 34 and an adder 35 for adding the output of the flip-flop 34 to the output of the loop filter 22. The numerical controlled oscillator 23 outputs digital oscillation frequency data according to a signal indicating the shift of the sampling timing output from the loop filter 22.\nThe tap coefficient arithmetic unit 25 provides a tap coefficient for advance or delay of a phase for the Fir filters 26a and 26b depending on the oscillation frequency data output from the numerical controlled oscillator 23, and outputs the I signal and the Q signal at the original sampling point from the Fir filters 26a and 26b. \nThe thinning control unit 24 controls the thinning circuits 27a and 27b, and thins data at an unnecessary sampling point in the data at each sampling point output from the Fir filters 26a and 26b. \nNext, the operation of the phase comparator 21 of the above-mentioned timing recovery loop is explained by referring to the operation explanation shown in FIG. 3.\nFIG. 3 shows the analog waveform of the input of the phase comparator 21. FIG. 3A shows the case in which the value of the input signal of the phase comparator 21 at each of the sampling time t\u22121, t and t+1 indicates a monotonous increase, and FIG. 3B shows the case of a monotonous decrease.\nIt is determined whether the values of the input signal indicate either a monotonous increase, based on the values of the input signal of the phase comparator 21 at each of the sampling time t\u22121, t and t+1, which are d(t\u22121), d(t), and d (t+1), a monotonous decrease, or any other cases, a predetermined arithmetic operation is performed on the value of an input signal at each sampling time based on the determination result, and outputs the arithmetic result as a phase determination result.\nThe phase comparator 21 outputs a value obtained by the following equation when the values of the adjacent input signal at the sampling time t\u22121, t, t+1 indicate a monotonous increase.\u22122\u00d7(d(t)+{d(t+1)\u2212d(t\u22121)}/2)\nThe phase comparator 21 outputs a value obtained by the following equation when the values of the input signal at the sampling time t\u22121, t, t+1 indicate a monotonous decrease.2\u00d7(d(t)+{d(t+1)\u2212d(t\u22121)}/2)\nFurthermore, the phase comparator 21 outputs \u201c0\u201d when the values of the input signal at the sampling time t\u22121, t, t+1 do not indicate a monotonous increase nor a monotonous decrease.\nFIG. 4A shows the output of the loop filter 22 when the frequency error between the original sampling point (hereinafter referred to as a symbol timing) and the sampling clock is small. FIG. 4B shows the output of the loop filter 22 when the frequency error against the symbol timing is large.\nWhen the frequency error is small, and the values of the input signal indicate a monotonous increase or a monotonous decrease, the phase comparator 21 outputs a value indicated in the above-mentioned equations, and the integrator 31 of the loop filter 22 integrates the values. As a result, for example, as shown in FIG. 4A, the output of the loop filter 22 increases until it reaches the convergence point of the timing recovery loop, and the capturing operation is completed at the convergence point.\nOn the other hand, when the frequency error is large, the signal points are distributed at random on the I and Q phase plane. Therefore, the average value of the output of the phase comparator 21 is approximately 0. As a result, the output of the loop filter 22 is fixed to a specific value, for example, as shown in FIG. 4B, close to \u201c0\u201d, or is fixed to a specific value.\nFIG. 5 shows the value of the limiter 32. When the average value of the output of the phase comparator 21 is approximately 0, the value of the limiter 32 to which an integration result is input is also a constant value.\nThe characteristic shown by the dotted line in FIG. 4B indicates the desired output characteristic of the loop filter 22 for reaching the convergence point, but the above-mentioned timing recovery circuit 19 cannot realize the capture characteristic as such.\nThat is, when the timing error is large, the output of the loop filter 22 is fixed to a specific value, and it is hard for the timing recovery circuit 19 to maintain the synchronization of clock timing.\nFIG. 6 shows the configuration of the carrier recovery loop for synchronization of the regenerated carrier wave with the carrier wave of a modulated signal.\nThe carrier recovery loop includes the rotor 17 and the carrier recovery circuit 18. The carrier recovery circuit 18 includes a phase comparator 41, a loop filter 42, a numerical controlled oscillator 43, and a Sin/Cos table 44.\nThe phase comparator 41 determines the advance or delay of the phase of the signal point on the I and Q phase plane obtained from the I signal and the Q signal to the phase of the normal signal point, and outputs a phase error signal for amendment of the advance or delay of the phase.\nThe loop filter 42 is a completely secondary loop filter, and includes a multiplier 45 constituting a low pass filter, a multiplier 46 constituting a high pass filter, an adder 47 for adding the output of the multiplier 45 to the output of an integrator 48, a limiter 49, the integrator 48, and an adder 50 for adding the output of the integrator 48 to the output of the multiplier 46. The integrator 48 is constituted by a flip-flop, etc. The coefficient \u03b1 provided for the multiplier 45 is a gain adjustment coefficient for a low pass filter, and the coefficient \u03b2 is a gain adjustment coefficient for the high pass filter.\nThe numerical controlled oscillator 43 includes a flip-flop 51 and an adder 52 for adding the output of the flip-flop 51 to the output of the loop filter 42. The numerical controlled oscillator 43 outputs oscillation frequency data for amendment of the frequency error (phase error) of a carrier.\nThe Sin/Cos table 44 is a table for generation of a sine wave and a cosine wave corresponding to the oscillation frequency data output from the numerical controlled oscillator 43.\nThe rotor 17 includes multipliers 53 and 54 for multiplying an I signal by the sine wave and the cosine wave output from the Sin/Cos table 44, multipliers 55 and 56 for multiplying a Q signal by the sine wave and the cosine wave, an adder 57 for adding the output of the multiplier 53 to the output of the multiplier 56, and an adder 58 for adding the output of the multiplier 54 to the output of the multiplier 55. (The output of the multiplier 54 is multiplied by \u201c\u22121\u201d and input to the adder 58.)\nThe rotor 17 multiplies the I signal and the Q signal by the sine wave and the cosine wave output from the Sin/Cos table 44 to rotate the signal point by the primary transform, and outputs the I signal and the Q signal at a desired sampling point.\nFIG. 7 shows a signal point of the I and Q phase plane (constellation point). The point indicated by the black dot shown in FIG. 7 is a normal signal point, the point indicated by a white dot is a signal point indicating a minus phase difference, and the point indicated by a dot with diagonal lines is a signal point having a plus phase difference. The counterclockwise rotation is a positive rotation and the clockwise rotation is a negative rotation.\nThe phase comparator 41 outputs a minus error signal when the position of the signal point on the I and Q phase plane determined according to the I signal and the Q signal output from the rotor 17 is detected in a predetermined range of the plus side relative to the position of the normal signal point, and outputs a plus error signal when it is detected in a predetermined range on the plus side on the I and Q phase plane. The phase comparator 41 outputs \u201c0\u201d when the position of the signal point depending on the I signal and the Q signal is detected in the position of the normal signal point.\nIn the carrier recovery loop, when the frequency error of a carrier is large, a signal point is distributed at random on the I and Q phase plane. Therefore, the average value of the output of the phase comparator 41 is approximately 0. Therefore, there has been the problem that, like the above-mentioned timing recovery loop, the output of the loop filter 42 is fixed to a specific value, and no carrier can be captured.\nThe patent document 1 describes providing a limiter for outputting an initialization value at the minimum level when the input value becomes larger than the limiter value in the feedback route to prevent pseudo synchronization.\nThe patent document 2 describes integrating a phase error signal, detecting synchronization by comparing an integrated value with a threshold, and stopping frequency sweep when synchronization is detected.\n[Patent Document 1] Japanese Patent No. 2885058\n[Patent Document 2] Japanese Patent Publication No. H5-30098"}
-{"text": "In a gas turbine engine fuel delivery system, pump assemblies, as shown for example in US 2005/0232784, are typically used for pumping the fuel. Where such assemblies include gear pump, gear elements are commonly supported by bearing blocks which are adapted to receive respective bearing shafts of the gears through a bore of each bearing block. These bearing blocks also typically abut axially-directed faces of respective gears of the pumps. The bearing blocks may be for solid bearings, or pressure loaded bearings. A solid bearing typically transfers load from journals to the pump housing, and additionally can transfer axial load to the housing. Pressure loaded bearings also transfer load from journals to housing, and in addition can provide an axial force and a moment against the axially-directed face of the gear which the bearing block abuts.\nIt is known to use bimetallic (alloy) bearing blocks, as shown, for example, in U.S. Pat. No. 4,523,365. Such a bimetallic bearing block generally comprises an inner bush covered with an outer backing layer. The inner bush is formed of an alloy which provides a tribologically compatible surface for the gear side face and journal to run against. However it may be undesirable for the whole bearing block to be formed of such an alloy. Two reasons for this are that firstly the weight of a block formed solely of such an alloy may be larger than desired, and secondly the difference between the coefficient of thermal expansion (CTE) and that of the light alloy normally used for the pump housing body may be large. Therefore, in a bimetallic bearing block the inner bush is coated with a light alloy backing layer, which reduces the overall weight of the block, and mitigates the CTE difference with the pump housing body.\nIn order to provide a tribologically compatible surface for the gear side face and journal to run against, an antifriction alloy is typically used to form the inner bush. The antifriction alloy may be, for example, a lead bronze alloy. In particular, the antifriction alloy may be a high concentration lead bronze alloy.\nHowever, testing of bearing blocks which use such an antifriction alloy has shown that such blocks can be prone to suffer permanent radial deformation along the bearing bore when operated above a threshold pressure/temperature combination. If this permanent radial deformation is significant, it can reduce the clearance between the gear journal and the bearing bore. As a result of this reduced clearance, overheating can occur, which may result in mechanical damage of the gear and/or bearing. This problem is more evident in bearing blocks having a higher concentration of lead than those having a lower concentration.\nFIG. 1 shows a typical Stribeck curve, which describes how coefficient of friction varies for different lubrication regimes. It can be seen that typically, higher coefficients of friction occur in a mixed mode lubrication regime as compared to in a full-film lubrication regime. The thrust faces in a gear pump bearing arrangement typically operate at least partly in a mixed-film lubrication regime. This is evident from the wear and scoring that can be visualised at the thrust surfaces after running. In contrast, the journal element of the bearing typically operates in a full-film lubrication regime. Generally, the loads on the thrust face of a bearing block are lower than those acting on the journal element of the bearing, but the sliding velocities are higher. Accordingly, local heating of the thrust face and gear side face is more likely than in the journal element of the bearing. The situation may be particularly acute at the region of the thrust face under the gear root circle diameter where fluid cooling is limited by the restricted fluid flow that occurs across this section of the thrust face. The primary mode for heat transfer away from the thrust face at this point may be conduction through the inner bush and backing layer of the bearing block. Therefore, depending on the thermal conductivity of the alloys of these components, there may be a problem that heat cannot transfer away from the thrust face sufficiently quickly. This may then lead to reduced wear characteristics, and correspondingly poorer performance of the thrust face in a thrust bearing capacity."}
-{"text": "1. Field of the Invention\nThe present invention relates to a component material for constructing a casing for various kinds of electronic devices.\n2. Description of the Related Art\nFor a casing for various kinds of electronic devices, such as a camera, a mobile phone, a PDA (Personal Digital Assistant), a personal computer, and so on, use has conventionally been made of materials such as metal (aluminum, stainless steel, titanium and magnesium) or a resin (acrylonitrile-butadiene-styrene (ABS) resin, polycarbonate, acrylic resin). These materials have been adopted not only from the functional consideration such as their moldability, rigidity and corrosion resistance, but also for aesthetics.\nRecently, further attention has also been paid to using, as a familiar, and environment-friendly material, wood or a woody material obtained by processed wood. For example, Jpn. Pat. Appln. KOKOKU Publication No. 7-24353 discloses an electromagnetic shielding material obtained through impregnation of wood or woody material with phenol resin and subsequent carbonization. The Publication discloses the idea of, while using the wood material, imparting an electromagnetic shielding property to a product obtained."}
-{"text": "The invention relates to a radial roller head according to the earlier part of patent claim 1.\nFrom DE 42 36 085 there is known a radial roller head with rollers arranged in a holder distributed at equal angular distances and at axial distances about the rolling axis. The roller surfaces have a spiral-shaped increasing course over the circumference opposed to their rotational direction. The rollers are coupled to a toothed gearing and a locking latches automatically after each complete roller revolution. Before the rolling procedure the locking is released, this being by the subject which is to be formed with the rollers. With the known radial roller head there is provided a stop which limits the introduction movement of the subject between the rollers, which is coupled to the locking of the rollers and which is displaceable in the direction of the rolling axis about an operating stroke against the action o)f compression springs. The locking is released on bearing of the subject due to its movement about the operating stroke.\nWith the known radial roller head there is the danger that during the releasing procedure or also thereafter a relative movement takes place between the subject and rollers if the subject has not achieved its stationary situation. This procedure naturally inhibits the forming procedure."}
-{"text": "The present invention is directed to a condensate drain pan for use under the heat exchange coil of an HVAC or refrigeration system. More specifically, the condensate drain pan is multi-functional in that it includes a bottom outlet and a side outlet and in that either outlet can be used on either end of an air conditioning cabinet.\nIn recent years, the air conditioning industry has had a heightened awareness of indoor air quality, particularly relating to the potential health concerns deriving from microbial growth in standing condensate. This standing condensate remains in the drain pans during the air conditioning unit's off cycle, and may be subsequently transported to an occupied space by way of the indoor air stream. For this reason, the air conditioning industry has recognized the need for a periodic cleaning of drain pans. Drain pans are positioned under a condensing heat exchange coil to collect the condensate as that condensate condenses and drips off the coil.\nTypical past industry practice has provided a flat condensate pan with condensate connections at only one end of the drain pan. In installing such drain pans, particularly in retrofit applications, it is not uncommon for the air conditioning unit to be improperly leveled during installation. If an air conditioning unit is improperly leveled such that the slope of the unit is away from the single drain pan connection, condensate will pool at the end of the drain pan opposite the drain connection. In more extreme cases, condensate will overflow the drain pan resulting in unwanted water in the base of the air conditioning unit and, potentially, in the occupied space and structure. Past solutions to this situation have included (a) removal of the air conditioning unit and re-leveling of the mounting curb at great expense, or (b) if the unit is large enough, providing a connection at either end of the drain pan. In many cases, the problem is simply ignored due to the difficulty in correcting the problem.\nFurthermore, access for cleaning of drain pans in the past has typically been through a time consuming removal of the exterior panels of the air conditioning units, followed by an attempt to clean the drain pan through a limited access aperture. As a result, the drain pan was often not cleaned, or was improperly cleaned. In some cases, chemical \"pills\" to counteract microbial growth have been added to drain pans in lieu of full drainage and proper cleaning.\nIndustry standards have been under development that make it unacceptable to provide air conditioning unit designs that allow condensate to pool in drain pans. The practical industry response to this has been to provide drain pans with sloped surfaces to ensure complete drainage. However, this assumes that the air conditioning units are properly leveled. With a sloped drain pan there is no longer a way to provide connections on both ends of the drain pan without one end of the drain pan being dysfunctional, thereby eliminating the second solution above and leaving only the solution of incurring the expense of re-leveling the air conditioning unit if the unit is improperly installed.\nRegarding bottom condensate connections, it is becoming common practice in the building/roofing industry to eliminate as many roof penetration apertures as possible due to the substantial warranties and potential liabilities associated with roof leaks in that industry. As such, there is a growing preference to bring the condensate disposal piping of the air conditioning unit from the space up through the bottom of the unit to eliminate a roof penetration and a potential leak related liability. Solutions to date have been (a) to route condensate piping out of the side of the unit and back in through the side of the rooftop curb, or (b) to provide both bottom and side outlet connections, requiring whichever outlet connection is not in use to be plugged or capped with a potential for leaks, or (c) to mount the drain pan in an elevated position within an air conditioning unit such that the piping could be routed out either the side or through the base of the unit.\nIt would be desirable to provide a drain pan which is easily cleanable, which is sloped, which can be installed from either side, and which includes both a bottom outlet and a side outlet."}
-{"text": "1. Technical Field\nThe present disclosure relates to monitoring systems and, particularly, to a monitoring system and a monitoring method.\n2. Description of the Related Art\nConventional monitor systems typically include a number of telecameras, a number of monitors, and a number of video cassette recorders. The telecameras aim at a scene so as to capture images of the scene from different directions, and output surveillance videos to the monitors for monitoring in real time. The surveillance videos are stored in the video cassette recorders.\nOne of the disadvantages of utilizing such a monitoring system is that in order to find a target person in the surveillance videos, all the surveillance videos need to be reloaded from the cassette recorders, and checked one by one. This manual searching task is time-consuming and inefficient.\nWhat is needed, therefore, is a monitoring system and a monitoring method, which can overcome the above-described problem."}
-{"text": "Charge-transporting thin-films composed of organic compounds are used as the emissive layer or charge injection layer in organic electroluminescence (organic EL) devices. Processes for forming such thin-films are broadly divided into dry processes such as vapor deposition and wet processes such as spin coating. In comparing dry processes and wet processes, wet processes are capable of efficiently producing thin-films having a high flatness over a large surface area. Therefore, in the organic EL field where thin-films of larger surface area are desired, organic thin-films are often formed by wet processes.\nIn recent years, there has been a steadily growing desire for higher functionality, greater multifunctionality, smaller size, etc. in manufactured products within the organic EL field. Under such circumstances, to achieve, for example, lighter weights and thinner dimensions, substrates composed of organic compounds have come to be used in place of glass substrates, and so the treatment temperatures that can be used in manufacturing processes have become lower than in the past. As a result, there exists today in the organic EL field an increasing desire for varnishes which can be baked at lower temperatures and which, moreover, even in such cases, provide thin-films endowed with good charge transportability."}
-{"text": "Many botanical substances contain chemicals that have been found to be useful for the therapeutic treatment of various medical conditions. Since these chemicals are often present in very small amounts, techniques have been developed to extract these substances and to concentrate the therapeutically active agents. Various methods are available for extraction and purification of such substances, including the use of organic solvents, microwave systems, and supercritical CO2 extraction. Organic solvent-based extractions utilize added solvents that are evaporated to form a concentrated extract, which results in a damp, pasty mass that is typically further spray-dried onto a carrier for delivery. Alternatively, supercritical CO2 extraction is another method of collecting such extracts. This extraction method yields a thick, high viscosity resin, oil, or other fluid-like material that can have a honey-like consistency.\nOne pharmaceutically useful botanical substance is the extract of hops (Humulus lupulus L.). Hops cone flowers contain a variety of active agents, including alpha acids, iso-alpha acids, and beta acids, as well as a number of flavonoids and essential oils. Humulone, one of the alpha acids found in hops, has been demonstrated to suppress cyclooxygenase-2 activity, inhibit angiogenesis, and decrease bone loss. Some other biologically relevant properties of hops constituents include anti-inflammatory, antibacterial, antiviral, antifungal, estrogenic, anti-oxidant, anti-allergenic, anti-carcinogenic, and anti-proliferative properties.\nAs with other botanical substances, dried hops flowers contain very small amounts of alpha acids. Supercritical CO2 extraction and other solvent-based extractions of dried hops cones produce a thick, high-viscosity resin that can contain a high percentage of active hops constituents. While extraction is an effective means of providing alpha and beta acids in a highly concentrated form, the resulting extracts have very low solubility in water. This property can make digesting such extracts difficult, resulting in delayed absorption of the acids and delayed onset of certain therapeutic effects. It would therefore be useful to provide the primary constituents of hops extracts in formulations that are soluble in water. In addition, methods of making such formulations from hops extract resin would be desirable."}
-{"text": "Fuel cells which generate electric current by controllably combining elemental hydrogen and oxygen are well known. One form of a fuel cell consists of an anodic layer, a cathodic layer, and a dense ion conducting electrolyte formed of a ceramic solid oxide. Such a fuel cell is known in the art as a \u201csolid-oxide fuel cell\u201d (SOFC). Hydrogen, either pure or reformed from hydrocarbons, is flowed along the outer surface of the anode and diffuses into the anode. Oxygen, typically from air, is flowed along the outer surface of the cathode and diffuses into the cathode. Each O2 molecule is split and reduced to two O\u22122 ions catalytically by the cathode. The oxygen ions are conducted through the electrolyte and combine at the anode/electrolyte interface with hydrogen ions to form molecules of water. The anode and the cathode are connected externally through the load to complete the circuit whereby electrons are transferred from the anode to the cathode. When hydrogen for the fuel cell is derived by \u201creforming\u201d hydrocarbons such as gasoline in the presence of limited oxygen, the \u201creformate\u201d gas includes CO which is converted to CO2 at the anode. Reformed gasoline and diesel oil are commonly used fuels in automotive fuel cell applications. However, other hydrogen-containing fuels for the reforming process such as, for example, JP8, natural gas, propane, synfuels, alkane alcohols, and coal based fuels may be used as well.\nA single cell is capable of generating a relatively small voltage and wattage, typically between about 0.5 volt and about 1.0 volt, depending upon load, and less than about 2 watts per cm2 of cell surface. Therefore, in practice it is known to stack together, in electrical series, a plurality of cells. The outermost interconnects of the stack define electric terminals, or \u201ccurrent collectors,\u201d which may be connected across a load. A typical prior art SOFC for use as an auxiliary power unit (APU) in a vehicle may comprise about 60 individual fuel cells and may generate, at full power, on the order of 5 kilowatts of electric power.\nA complete SOFC system typically includes auxiliary subsystems for, among other requirements, generating fuel by reforming hydrocarbons as discussed above; tempering the reformate fuel and air entering the stack; providing air to the hydrocarbon reformer; providing air to the cathodes for reaction with hydrogen in the fuel cell stack; and providing air for cooling the fuel cell stack.\nA known shortcoming of a complete SOFC system is that it inherently has a relatively large thermal mass, and consequently, such a system is relatively slow in ramping up to full electric output. Electric output can't begin until the fuel cells are warmed to about 550\u00b0 C., and a temperature of about 750\u00b0 C. is required for full output and steady-state operation. The mass of the fuel cell stack along with the induced thermal stresses caused by heating the stack through heated air applied to the cathode dictate the length of time required to produce electricity. Even with known methods for preheating and forced heating of elements in an SOFC system, a prior art SOFC system requires on the order of sixty minutes, starting from ambient temperature, to begin producing usable amounts of electricity. Thus, costly and bulky additional energy storage systems would be required on an APU-powered vehicle to provide power during the fuel cell warm-up period, which causes dissatisfaction to users of such systems. In addition, long start-up times result in a reduction in operating efficiency, especially for intermittent, short-duration operation.\nDuring these start-up periods, unwanted reactions of reformate fuel in the stack is possible, causing deposition of soot (coke) on the relatively cool anode surfaces. Such deposition is undesirable and can result in degradation and eventual failure of the fuel cell stack.\nSteam is a known preventor of carbon (coke soot) forming reactions and a known cleaner of soot from anodes. Therefore, the anodes may be cleaned and prevented from coking by injection of hot steam into the fuel cell stack during warm-up to an operating temperature at which coking does not occur. One approach is to produce steam via vaporizing water from a separate water storage tank, and to then inject the steam into the fuel cell stack along with the reformate. Such a process is undesirable since it requires added complexity of apparatus and logic for storing, supplying, vaporizing, and replenishing water adjacent to the fuel cell system.\nWhat is needed is a means for reducing the start-up period required to bring a large solid-oxide fuel cell system to operating temperature.\nWhat is further needed is a means for preventing coke from depositing on the anodes of a solid-oxide fuel cell during the start-up period required to bring the fuel cell system to operating temperature.\nWhat is still further needed is a means for providing hot steam to an SOFC fell cell stack to remove coke which may form on the anodes while the anodes are cool.\nIt is a principal object of the present invention to reduce the start-up time required for a solid oxide fuel cell system to begin producing electricity.\nIt is a further object of the invention to reduce the need for electrical storage systems in applications wherein a solid oxide fuel cell is an auxiliary electric power unit.\nIt is a still further object of the present invention to minimize coking of anodes during the start-up period required to bring a fuel cell system to operating temperature.\nIt is a still further object of the invention to remove coke deposits that have accumulated on the SOFC anodes."}
-{"text": "Sleep is essential for health and quality of life. Insomnia is a growing health problem in the United States. It is believed that more than 10-15 million people suffer from chronic insomnia and up to an additional 70 million people suffer from some form of insomnia each year. Insomnia is a condition characterized by difficulty falling asleep (sleep onset), waking frequently during the night (fragmented sleep), waking too early (premature final awakening), and/or waking up feeling un-refreshed. In the National Sleep Foundation's (NSF) Sleep in America Poll 2005, 42% of survey respondents reported that they awoke frequently during the night, 22% of adults reported waking too early and not being able to return to sleep and 38% reported waking and feeling un-refreshed.\nSleep maintenance difficulty is the most commonly reported symptom in primary care patients with chronic insomnia, and is the most common complaint in depressed patients, medically ill populations, especially those with pain symptoms, and in the elderly.\nMedications commonly used to treat sleep disorders, such as insomnia, include sedative antidepressants, antihistamines, benzodiazepines, and non-benzodiazepine hypnotics.\nAlthough there have been several advances in pharmaceutical treatments for insomnia, it is often hard to find an ideal drug for treating particular forms of insomnia. One common problem is early termination of sleep or premature final awakening. For example, many individuals may wake prematurely and not fall back asleep, thereby failing to achieve a full night of sleep. Many drugs that are effective in inducing or expediting sleep initiation do not provide much effect in maintaining sleep, particularly through the eighth and final hour of sleep period. Drugs that are sufficiently powerful to induce a full eight hours sleep often cause serious hangover effects, i.e., the patient has difficulty awakening and/or feels sedated, sleepy, or disoriented and may demonstrate impairment of psychomotor function.\nIn addition to patients having difficulty with early termination of sleep during the last 60, 90, or 120 minutes of an 8 hour sleep period, other patients have problems with fragmented or disrupted sleep. In other words, those patients awaken one or more times during that time period, then fall asleep again. Such fragmented sleep patterns detract from a feeling of restfulness, and make it less likely that the patient will enjoy restful sleep.\nBoth groups of patients would benefit greatly from a drug that addresses their particular sleep deficiency.\nDoxepin is a tricyclic antidepressant that is known to have beneficial effects in treating insomnia. See, e.g., U.S. Pat. Nos. 5,502,047 and 6,211,229. However, prior to the present invention, doxepin was not known to have particular efficacy in treating premature termination of sleep at the end of an 8 hour sleep period, nor was it known to be efficacious in treating those patients with disturbed sleep patterns during the final 60, 90, or 120 minutes of an 8-hour sleep period. The mean half-life of doxepin is 17 hours, and the half-life of its major active metabolite, desmethyldoxepin, is 51 hours. Thus, when taken at the start of a sleep cycle, a majority of the drug or active metabolite should still be present in the body at the end of the sleep cycle. As a result, it would be expected that dosages of doxepin that are sufficient to address premature final awakenings or last-hour sleep efficiency in the elderly would also cause post-sleep sedation or other undesirable side effects.\nThe present invention describes the surprising ability of doxepin to treat last-hour sleep efficiency and premature final awakenings in patients, without untoward side effects."}
-{"text": "Devices within a distributed computing environment must typically exchange large amounts of configuration data in order to identify and access various distributed resources. For example, a central server within a client-server computing model may be required to distribute a large amount of configuration data to each client within its distributed computing environment in order to enable these clients to identify and access the distributed resources they require to perform various tasks. Similarly, a peer within a peer-to-peer computing model may request a large amount of configuration data from one or more of its peers in order to identify and access the distributed resources the peer requires to perform various tasks.\nA distributed computing system's configuration data is typically represented in the form of a tree that identifies the hierarchical relationship of each object within the distributed computing system. Nodes within this tree may represent objects (such as a cluster, system, group, application, resource, or attribute of the same) within the distributed computing system connected in parent-child relationships. Devices within the distributed computing system may be interested in either the entire tree (e.g., peers within a peer-to-peer model or a server within a client-server model may be interested in the entire tree) or a subset of the tree (e.g., clients within a client-server model may only be interested in a subset of the tree).\nUnfortunately, due to the typical size of this configuration data, an inordinate amount of network bandwidth may be required to exchange even a subset of this configuration data among devices within the distributed computing system, resulting in potential delays. For example, the amount of time required for a client to recover from a failure may be dramatically increased due to the time required to request, receive, and process a snapshot of (even a subset of) a central server's configuration data. As such, the instant disclosure identifies a need for systems and methods for efficiently synchronizing configuration data within distributed computing systems."}
-{"text": "My U.S. Pat. No. 4,367,850 describes a cartridge including a length of magnetic tape spliced into an endless loop. In one embodiment that cartridge can be made to conform to the standard of the NAB so that it can be used in the types of record/playback machines presently in use in the broadcast industry. That cartridge embodiment causes much less edge wear on the tape and much less change in tension in the tape at a head of a machine in which the cartridge is received than the other types of cartridges presently in use in the broadcast industry.\nThe cartridge described in U.S. Pat. No. 4,367,850 employs a method for maintaining a uniform high tension in the tape at a head of a machine in which the cartridge is engaged so that the tape can be pressed against the head by tension alone without the use of pressure pads. That method comprises providing a fixed, generally cylindrical hub having a central opening and a slot extending axially across the full width of the hub and communicating with its central opening, and an endless length of magnetic tape; wrapping a major portion of the tape about the hub to form a coil while allowing a minor portion of the tape to extend from the innermost wrap of the coil through the slot into the central opening of the hub, and around the side surface of the coil to the outermost wrap of the coil; pulling the tape from the slot and across a head on a record/playback machine; and applying a light force to tension the minor portion of the tape as it moves onto the outer wrap of the coil which will produce the high, generally uniform tension in the minor portion of the tape being pulled from the coil and across the head.\nThe embodiment of the cartridge adapted for use in the broadcast industry which employs that method to tension tape across a head comprises a housing adapted to be received in a record/playback machine and having access openings adapted to receive record/playback heads and a tape drive mechanism in the machine. The hub is fixed on the housing at a position spaced from the access openings. Means on the housing define a tape path for, and produce tension in, the minor portion of the tape. Those means comprise means for guiding the minor portion of the tape past the access openings in a predetermined position for engagement by the heads and the drive mechanism of the playback machine, and a guide pin guiding the minor portion of the tape between the access openings and the outermost wrap of the coil. The guide pin is mounted for movement between a first position providing a first path length between the access openings and the outer wrap of the coil and a second position providing a second path length longer than the first path length between the access openings and the outer wrap of the coil (which means preferably is an arm having a first end supporting the guide pin, and a second end pivotably mounted on the housing to afford movement of the pin along an arcuate path adjacent the periphery of the coil between its first and second positions), and means are provided for biasing the guide pin toward its second position.\nThe pin can move to positions between its first and second positions under the influence of the biasing means to accommodate changes in length of the minor portion of the tape which decreases or increases respectively when the length of the major portion increases or decreases. The major portion of the tape cyclically undergoes its largest change in length by slowly decreasing in length as the splice moves from the coils outer wrap toward its inner wrap and by then suddenly increasing in length as the splice leaves the coil. The means for biasing the guide pin is adapted to apply a small force at the guide pin (e.g., generally in the range of 2 to 14 grams) to lightly tension the tape extending around the pin and moving onto the outermost wrap of the coil, somewhat in the manner of a rope or Proney brake, which light tension produces a significantly higher tension (e.g., generally in the range of 50 to 115 grams or 2 to 4 ounces) in tape leaving the coil. That higher tension is surprisingly uniform despite small changes in the force applied by the guide pin as the length of the minor portion of the tape changes due to the position of the splice along the tape.\nThe tape is guided so that the quite uniform higher tension thus produced in the minor portion of the tape between the inner wrap of the coil and the drive mechanism of a machine in which the cartridge is engaged presses the tape against the record and playback heads of the machine with sufficient pressure that pressure pads are not required, and the tape is thus not subjected to the abrasion and erratic forces caused by sliding contact between the tape and such pressure pads.\nAlso, my U.S. Pat. No. 4,394,989 describes a tape cartridge of the type described above in which the length of the minor portion of the tape can be easily adjusted by persons in the field to accommodate changes in the cartridge such as tape wear to ensure that a desired range of guide pin movement will be retained.\nIn that improved cartridge, the hub comprises a flexible arcuate cantilevered portion having a first end partially defining the slot, a second end that is fixed on the housing, and which is separated from the housing between its ends. Means in the form of a cam rotatably mounted on the housing is provided for changing the position of the cantilevered hub portion radially of the rest of the hub to adjust the effective diameter of the hub and thereby the length of the minor portion of the tape to cause the guide pin to move within a desired range of movement between its first and second positions. The cam has a peripheral surface contacting the inner side of the cantilevered hub portion, and can be manually rotated with a tool such as a screwdriver as the tape is being propelled by a machine so that the effects of cam adjustments on the tape can be observed as they are made. Such adjustability ensures that changes in overall tape length which can occur due to causes such as tape wear can be compensated for. While such adjustments do not need to be made often during the life of a cartridge, they do require attention by technicians, such as each time the tape is re-recorded.\nThus the cartridge described in my U.S. Pat. No. 4,394,989 is not particularly useful in cartridges used in background or foreground music machines where users typically are not trained to or interested in making the occasional cam adjustments needed to maintain the cartridge in its best operating condition.\nAnother problem that occurs with the cartridge described in U.S. Pat. No. 4,394,989 is that when it is moved from a room at normal temperatures into the cold (e.g., from an area at about 68.degree. F. to an area at below 10.degree. F.) its coil and fixed hub will shrink. When the cartridge is then again brought into a warm room and playing of the cartridge is attempted, its tape coil will seize around the hub which is more exposed and returns to room temperature and thus expands to its original size more quickly than the tape coil. Radio technicians can correct for this problem by adjusting the cam to release the tension in the coil. Again, however, persons involved with background music or foreground machines typically lack the training and desire to perform similar adjustments."}
-{"text": "Modern large commercial buildings, such as factories, hotels, office buildings and hospitals, frequently use large and complex heating, ventilating and air conditioning (HVAC) equipment.\nIt is known to equip commercial buildings with variable air volume systems, which are capable of meeting the entire cooling and heating requirements of the building. Within the building, there are likely to be located a number of terminal units in different zones throughout the building, each connected via duct work to a central air supply. Such terminal units are sized to meet the conditions of the space which each serves, but, as a result, multiple offices, rooms or compartments within the structure are necessarily supplied with heating and cooling air by one terminal unit.\nThe end result of this type of design is that individual rooms or compartments within a structure are forced to share a common heating and cooling environment. While this may represent nothing more than a minor inconvenience for many building occupants in most cases, it presents particular difficulties in some specific environments, for example, hospital operating rooms.\nPrecise control of temperature and humidity in hospital operating rooms is important. Such rooms are frequently equipped with a number of machines which generate substantial heat. Further, the rooms will be populated with a varying number of workers during a typical operative procedure. Further, operating rooms must be regularly reconfigured for different procedures, meaning that the equipment and personnel contained within the room will vary substantially from day to day.\nUnder these circumstances, it is extremely difficult to maintain desired, consistent temperature and humidity levels in specific areas within buildings where centralized heating, ventilating and air conditioning systems are in use.\nWhile it is known to install modular heating, ventilating and air conditioning systems in individual rooms and compartments, many such devices are inefficient, cumbersome to install, and take up substantial space in the room in which they are installed. Further, such stand alone units are not centrally located within the rooms or compartments which they are designed to service, resulting in an imbalance in temperature and humidity in different areas of the same room or compartment. Further, such self-contained units often recirculate, rather than vent room air. Such units sometimes are in conflict in operation with the building central heating, ventilating and air conditioning system, resulting in energy inefficiencies when a local modular unit attempts to heat the air within a particular room or compartment at the same time as the centralized heating, ventilating and air conditioning system is attempting to cool the very same space.\nIt is desirable, therefore, to implement a modular heating, cooling and humidifying system which works in concert with the centralized heating, ventilating and air conditioning system of a larger structure, and which can be placed within the air ducting system of an existing structure, allowing individual temperature and humidity control in a single compartment or room, while at the same time not occupying physical space within the room or compartment, and further operating in symbiosis with the central heating, ventilating and air conditioning system of the structure."}
-{"text": "Malaria has a tremendous impact on human health, killing millions annually and the disease is a major impediment for social and economic development of nations in malaria-endemic areas, particularly in sub-Saharan Africa (Sachs & Malaney (2002) Nature 415:680-85). Malaria is a mosquito-borne disease that is transmitted by inoculation of the Plasmodium parasite sporozoite stage. Sporozoites invade hepatocytes (Kappe et al. (2003) Trends Parasitol. 19:135-43), transform into liver stages, and subsequent liver stage development ultimately results in release of pathogenic merozoites that initiate the blood stage cycle of infection (Shortt & Garnham (1948) Nature 161:126).\nProtection against blood stage malaria can be passively transferred by antibodies. Effectiveness of passive transfer of IgG between adults and children living in different geographic regions indicate that the antigens that are targeted by antibodies that protect against blood stage malaria are conserved (see Duffy et al. (2005) Vaccine 23 (17-18):2235-42). The best evidence that naturally occurring immunity against blood stage malaria targets the IRBC surface comes from studies of pregnancy malaria. In areas of stable malaria transmission, adults enjoy immunity that limits parasitemia and prevents disease. Women become more susceptible during pregnancy, and previously this was ascribed to pregnancy-related immunomodulation that develops to prevent rejection of the fetal allograft. However, susceptibility is greatest in primigravid women and diminishes over successive pregnancies, suggesting an acquired immune response to an antigenically distinct microbe. Placental isolates of P. falciparum commonly bind to chondroitin sulfate A (CSA) expressed on the surface of the syncytiotrophoblast (the cells lining the maternal vascular space in the placenta), and this phenotype is uncommon among isolates obtained from non-pregnant individuals (Fried & Duffy (1996) Science 272:1502-4). The distinct binding phenotype renders pregnant women na\u00efve to this parasite population during their first pregnancy (primigravidas). Women with multiple pregnancies (multigravidas) residing in areas of stable malaria transmission develop antibodies that inhibit parasite adhesion to CSA (Fried et al. (1998) Nature 395:851-2). These anti-adhesion antibodies are associated with a reduced mass of parasites in the placenta, and substantial improvements in fetal development (Duffy & Fried (2003) Infect. Immun. 71:6620-3). Because CSA-binding parasites cross-react with sera donated by mothers throughout Asia and Africa, the antigen targeted by protective antibodies is presumed to have conserved features or limited diversity.\nThere is a need in the art for vaccines that protect against malaria infection and disease. There is also a need in the art for diagnostic markers for malaria. The present invention addresses these needs and others."}
-{"text": "For the purposes of operation and in order to store information, computer systems and mainframes have a memory arrangement comprising a multiplicity of memory modules which are arranged, for example, in memory cabinets and have semiconductor components grouped on the memory modules.\nSemiconductor components, for example, DRAMs (Dynamic Random Access Memories), are generally subject to extensive function tests early in the production process before final assembly to form a memory module. These function tests are used to identify faulty memory cells or faulty column lines or row lines or generally faulty circuit parts in the semiconductor components. To this end, data values are written to memory cells in a memory cell array in the semiconductor component and are then read out again in order to be compared with the prescribed data values. This makes it possible to test the semiconductor components under various operating conditions in order to guarantee fault-free operation of the memory chip.\nThe memory modules accommodate a test device, a \u201cBIST\u201d (Built-In Self-Test) unit, in each of the semiconductor components as part of the latter. The BIST unit integrated in the respective semiconductor component carries out the requisite electrical function tests before the semiconductor components are installed. The BIST unit has a BIST controller which, as a switching region in the semiconductor component, is in the form of an ASIC (Application-Specific Integrated Circuit). Commands in a test sequence which are issued by the BIST controller are forwarded to the semiconductor component, with the BIST controller monitoring and evaluating the execution of the commands. The data transmitted by the semiconductor component regarding its operating states are output, for example, to external test systems which make an appropriate evaluation on the basis of which it is possible to make a statement regarding whether and, if appropriate, which memory areas are not functioning as intended. When the tests are carried out successfully, the module is classified as functional and is used in the customer's target application.\nHowever, a meaningful test result can only be achieved when the semiconductor component is tested at the operating frequency that it has during normal operation. A fault in a semiconductor component is always associated with the target application, for example, a voltage supply or an input parameter for configuring the semiconductor component is not within the prescribed specification.\nHowever, today's test methods do not yet make it possible to simulate these or all characteristic operating modes of the application in order to test the semiconductor components in proximity to the application. It is thus not possible to make a statement regarding whether the semiconductor components tested during production will run through all of the operating modes, occurring during later application without any faults.\nIn the event of a fault occurring during normal operation at the customer's premises, the memory module has to be returned to the manufacturer for the purposes of analysis. Identification data which may have been programmed in, for example the chip ID, test data or adjustment parameters, can thus be used only for subsequent historical tracking, but not for user-specific adjustment during normal operation.\nAt such a point in time it would be desirable, for the purposes of evaluation and/or analysis, to use a test system which would make it possible to test and adjust the semiconductor components during normal operation. Today's available external test systems are connected to the semiconductor components on a memory module via the semiconductor component's standard interface, which is used for external data interchange, address interchange and/or command interchange during normal operation. In the test mode, the test system is able to generate the test commands required to test the memory module, such as control and address commands, commands for reading and storing data words and also a clock signal, and is able to initiate the electrical function test via the BIST unit, for example. In the case of semiconductor components which can be operated in parallel, however, this function test may usually be carried out only for all chips, that is, all of the semiconductor components arranged on the memory module are tested in parallel at the same time via the standard interface.\nThe desire for remote access monitoring and a fault analysis capability for computer systems, i.e., driving the computer systems via an external test system, has existed since computer systems were networked. One difficulty in checking individual hardware components in a computer system is based on the fact that respective redundancy must be available for faulty components.\nFor example, an apparatus for the remote monitoring of computer components in a computer system is known, but this apparatus explicitly requires a working memory in the computer system. In order to analyze a processor in a computer system, a remotely accessible integrated debug environment has been proposed, wherein a computer which is connected to the computer system via the Internet can analyze the processor in the event of a fault. Further, a method for remotely accessing a faulty booting computer is known, in which the computer, in the event of a failed starting attempt, has recourse to a simple E-BIOS code which connects the computer to a service computer via a LAN or an Internet connection and thus makes the computer accessible for remote access operations for the purpose of repair and/or diagnosis. A known method for testing an SDRAM in a computer system, including test modules integrated in the computer system, tests the memories before starting by appropriate test modes in order to boost or attenuate a possible fault mode.\nIn addition, there is known a semiconductor module having semiconductor components which are arranged on the semiconductor module and are connected to one another via a serial line and having an interface which accesses the semiconductor components via the serial line.\nHowever, remote access maintenance of memory modules during normal operation has not been disclosed. To this end, it would be necessary to operate the computer system normally and to carry out specific addressing operations during an application.\nA memory arrangement which makes it possible to drive the semiconductor components (arranged on the memory arrangement) during normal operation using remote access and to test and adjust them in proximity to the application without impairing the operation of the semiconductor components is desirable."}
-{"text": "The present invention relates to a method for rotating at least two-dimensional image records with non-isotropic topical resolution.\nKnown literature contains a number of methods for rotating and shifting image records. These methods can be roughly divided into two categories. In the first category, the rotations and translations of a record are computed exclusively in the image space. In the second, the rotations and translations are computed completely or partially in Fourier space with the aid of Fourier transformation.\nThe methods of the first category are disclosed in the following examples: Starting with an existing image record, a rotated image record is obtained by executing the following two steps: In the first step, the coordinate points of the rotated image record are computed by multiplying the coordinates of the existing image record by a rotation matrix. The value of each rotated coordinate pointxe2x80x94this can be, for example, a matter of the percent value on a gray scalexe2x80x94is calculated in an interpolation process. A large number, more or less, of values of neighboring points of the existing image record, or all values, should be integrated into the interpolation process. The interpolations can be based on a Sinc interpolation, for instance. The following example should demonstrate that said interpolation processes entail a significant computing outlay. Assume a two-dimensional image with 256 pixels in two dimensions. This means that 65536 pixels must be rotated. In turn, this means that in the worst case for each of the 65536 pixels a respective interpolation must be performed with the same number of pixels. Altogether, this means a quite significant computing outlay and thus time outlay. For shifts of image records in accordance with methods of the first category, the steps described above for a rotation are performed accordingly.\nThe methods of the second category are exemplified next. There, Fourier transformation is specifically used to reduce the computing outlay. The Fourier transform F(k) of a one-dimensional image space function f(x) is generally defined as       F    \u2061          (      k      )        =            \u222b              -        \u221e                    +        \u221e              \u2062                            f          \u2061                      (            x            )                          \u00b7                  \u2147                      j            \u00b7            2            \u00b7            \u03c0            \u00b7            k            \u00b7            x                              \u2062                        \u2146          x                .            \nFor the inverse-transformation,       f    \u2061          (      x      )        =            \u222b              -        \u221e                    +        \u221e              \u2062                            F          \u2061                      (            k            )                          \u00b7                  \u2147                                    -              j                        \u00b7            2            \u00b7            \u03c0            \u00b7            k            \u00b7            x                              \u2062                        \u2146          k                .            \nAs the central feature of the Fourier transformation, the displacement set is of primary importance in the methods of second category. The displacement set for a one-dimensional Fourier transformation is as follows:\nf(x-x0) in the image space corresponds to ejxc2x72xc2x7xcfx80xc2x7x0xc2x7kxc2x7F(k) in Fourier space.\nApplied to multidimensional applications, this means that a shift of an image record by a random vector is imaged in Fourier space in an additional phase of the Fourier transform. Unlike methods of the first category, interpolations are not required for the translation of an image record. The interpolation is completed on the basis of the attributes of Fourier transformation semi-automatically with the transformation into Fourier space, with the introduction of a desired additional phase and the corresponding back-transformation. Rotations of image records can likewise be performed with the aid of Fourier transformation and its displacement set without interpolation. To accomplish this, a rotation matrix describing the rotation is dismantled into a product of corresponding shear matrices. The effect of shear matrices is expressed in simple shifts of rows or columns of an image record, respectively. These shifts can be calculated easily in Fourier space with the aid of the Fourier transformation and its displacement set. A method of another kind is described in the essay by William F. Eddy, Mark Fitzgerald and Douglas C. Noll, xe2x80x9cImproved Image Registration by Using Fourier Interpolationxe2x80x9d, MRM 36, pp. 923-931, 1996.\nThe case of a non-isotropic topical resolution, which is the more frequent case in many applications, can be handled by the above method of the second category only with additional outlay. Non-isotropic records are not problematic for methods of the first category. For example, in rotations an elliptical path is merely executed instead of a circular path. The high computing outlay for the necessary interpolations remains unchanged. In methods of the second category, the non-isotropic record must be converted into an isotropic auxiliary record in a first step. Then, a method of the second category can be applied. Finally, the record must be converted back into a non-isotropic record. These conversion processes, which are known as resampling processes, represent additional computing outlay, regardless of the type and manner of execution of the resampling processes. In practice, the registering of non-isotropic records is often avoided for these reasons.\nIt is thus an object of the invention to create a method which reduces the computing outlay in rotations and translations of non-isotropic image records.\nThis object is inventively achieved in accordance with the present invention in a method for rotating at least two-dimensional image records with non-isotropic topical resolution, said method comprising the steps of: describing a rotation of an image record using a rotation matrix; representing said rotation matrix as a product of at least two shear matrices, each shear matrix having at least one element that is dependent on the angle of said rotation and remaining elements which are exclusively zeroes and ones; multiplying a matrix element, which is dependent on said angle of said rotation, of at least one of said shear matrices by a factor; and performing said rotation of said image record in Fourier space without interpolations and without forming an isotropic auxiliary record, upon exploitation of a displacement set of the Fourier transformation by implementing said shear matrices as displacements of line elements of said image record.\nIn an embodiment, the present method provides the following steps:\ndescribing a rotation of an image record with a rotation matrix;\nrepresenting the rotation matrix as a product of at least two shear matrices, each of which comprises exactly one element that is dependent on the angle of rotation, and whose remaining elements are exclusively zeroes and ones;\nmultiplying the matrix element, which depends on the angle of rotation, of at least one shear matrix by a factor; and\nperforming the rotation of the image record in Fourier space, without interpolations and without forming an isotropic auxiliary record, by exploiting the displacement set of the Fourier transformation by implementing the shear matrices as displacements of line elements of the image record.\nBy simple one-time multiplications, by factors, of matrix elements of shear matrices that describe rotation intensive transformations of non-isotropic image records into isotropic auxiliary image records, corresponding back-transformations are made superfluous for the purpose of rotating non-isotropic image records in Fourier space. The application of the invention in a computer system reduces the necessary computing time.\nIn an advantageous embodiment, the factor depends exclusively on at least one ratio of a topical resolution in a first dimension to the topical resolution in a second dimension, which ratio describes a non-isotropic topical resolution. With the aid of a factor that is selected in this way, the error between the image record that is rotated as desired and the image record that is rotated by applying the method is minimized.\nIn a particularly advantageous embodiment, said error equals zero for two- or three-dimensional image records.\nAn advantageous embodiment relates to image records that are generated by a computer or MR tomography device. Here, the present invention can be advantageously employed particularly in the field of functional MR tomography, where the image records that are registered in a time series are rotated and displaced continuously for the purpose of motion correction."}
-{"text": "1. the Field of the Invention\nThis invention relates to sound generation and, more particularly, to novel systems and methods for providing aural and tactile signals to a human being.\n2. The Background Art\nAutism is a condition characterized by persons being overwhelmed by sensory perceptions, which become distorted (disoriented). Typically, the life of an autistic individual is characterized by overwhelming chaos due to sensitivity to external stimuli, including sight, sound, touch, taste, and smell. Various systems have been developed for treating autism.\nPrevious work has been done by the inventor in the area of orientation counseling and procedures. It has been provided in the past to place a headset or earbuds on a user in order to generate a sound that is detected by the ears, processed by the brain, and treated as a focal point or an orientation point for an individual. Once focused, perceptions become accurate or undistorted (oriented). However, many persons, particularly young children or severely affected subjects, may not be able to wear a headset in order to listen to the sounds that will help to focus attention, direct the \u201cmind's eye,\u201d toward the sound perceived to be at a particular position in space.\nFor example, the essence of stereo is the stereophonic projection of sound from different locations. Stereophonic headphones, for example, permit different tracks of music or other sound to be passed to the ears, thus giving an impression of position and distance. In reality, stereo sound is propagated by a speaker relatively nearer each ear. However, by changing the volume of different tracks, particular sounds may be played back to a listener in such a way as to appear to be originating at different locations, to move with time, or the like. However, experts on the subject have informed the inventor that it is impossible to exactly locate in space a perceived location for the sound. More particularly, the inventor was informed that it was impossible to maintain such a location with respect to an individual.\nThe David Autism Approach is documented in various books and websites available internationally. Unfortunately, many autistic individuals are unable to wear headphones, or unwilling to do so. Thus, an ability to receive a sound signal, which has been generated with certain artifacts (such as accommodating the geometry and attenuation of the hearer's head) and balanced to seem like it is coming from a particular point in space fixed with respect to the person, is not heretofore possible. However, it would be an advance in the art to provide a system, that generates sound, in a way that such a system can be created and then worn in such a location and manner that it is not easily removed by a subject, and does not interfere with movement of the head, and normal daily activities."}
-{"text": "Characteristic of cardiothoracic surgery is the post-operative patient who is sent to the Intensive Care Unit (ICU) intubated due to respiratory requirements. Approximately half of these patients are extubated within their first twenty-four post-operative hours. In most cases these patients are extubated within the first three days. There are some, however, who remain intubated for a significant length of time. When a surgeon identifies a patient who requires intubation longer than seven days, the surgeon will usually decide to perform a tracheotomy on that patient. The breathing support tube enters the trachea rather than entering the mouth for the trached patient. Communication for a intubated or trached patient is minimal due to the inability to speak resulting in the patient, hospital staff and loved ones resorting to the reading of lips, nodding of heads and squeezing of hands to communicate.\nWithout effective communication, the intubated or trached patient may not receive the standard of care he or she would otherwise receive had he or she been able to effectively communicate. The lack of communication also creates unnecessary levels of anxiety which the patient must endure. Nurses and hospital staff ask many questions from the patient pertaining to their prognosis and progress which may never get fully or even adequately answered. A doctor or nurse is not able to treat a symptom which they know little or nothing about. In addition, other problems arise due to the insufficient communication from the patient. Localized areas of pain are often mis-diagnosed, resulting in over-medication generally or the medication of an area which is not the source of pain. Proper and essential treatment given in an adequate and timely manner will help resolve or prevent many post-operative complications and decrease the patient's length of stay in the hospital. This begins with providing the patient a clear and precise means of communication.\nAnother problem exists in that currently patients are subjected to pushing a button or call light, which turns on a light in the hallway at the doorway to their room. Nurses have no way of identifying whether the patient's need is urgent or non-urgent. Additionally, the nurse is unable to prepare him/herself for the need appropriately before entering the room. Instead, the nurse must go to the patient's room, be informed of the problem or need and then leave the patient's room and retrieve whatever resources are necessary for the nurse to fulfill the patient's need or request. This is extremely time-consuming, wastes precious hospital resources, and can delay meeting patient's needs. This problem can be detrimental to the patient when the need is of an urgent matter. Unless the patient can scream loud enough to be heard from wherever help may be, the patient is subjected to wait until someone responds to a common light at the patient's doorway.\nMoreover, current systems lack a patient-centric device for the bedside interface (e.g., pillow speaker); and, only a few nurse call systems provide an opportunity for patients to convey specific messages directly to their assigned providers.\nFurthermore, eighty percent of hospitals care for patients with limited English proficiency (LEP) on a regular basis, and despite advancements in the profession of healthcare interpreting and translation services (aka, Language Access Services), patients with language barriers are often left without an effective means to communicate with their providers. While best practice, clinical ethics and legal and regulatory guidelines recommend the use of professional interpreters for all healthcare encounters, logistics and resource capacity make this prohibitive. Reasons cited by hospital staff for not using professional interpreters include resources available to bridge the language barrier are not user friendly; resources are not easily accessible; staff are unaware of the resource and have not been trained. Despite these disadvantages, nurse call systems do not provide a means for LEP patients to generate a nurse call request in the patient's preferred language.\nAccordingly, it has been estimated that inefficient communication costs U.S. hospitals more than $12 billion annually or $4 million for each 500 bed hospital. In summary, nurse call systems have been the primary means for hospitalized patients to initiate an encounter from the bedside. However, these nurse call patient requests range in urgency, are not differentiated based on skill-set required to fulfill the patient request, and are not equitable for LEP patients. Further, these shortcomings prohibit effective communication with LEP patients and can contribute to poorer outcomes compared to their English-speaking counterparts."}
-{"text": "A Fabry-Perot interferometer is an instrument commonly used in high resolution optical spectroscopy, and as a tool for the analysis of laser radiation. Etalons are widely used in telecommunications, lasers, and spectroscopy for controlling and measuring the wavelength of light. For example, a laser beam typically comprises several discrete optical frequencies related to different modes of oscillation of the laser resonator. The Fabry-Perot interferometer can be used to determine the number of modes oscillating.\nThe Fabry-Perot interferometer comprises two parallel mirrors and a Fabry-Perot etalon which has certain optical resonator properties. One of the mirrors is fabricated at different distances (two or more levels) from the other mirror. Conventional spectrometers comprise a large, expensive, and complex system of lenses and mirrors. More recently, attempts have been made to manufacture a Fabry-Perot interferometer within a monolithically integrated circuit. See for example, \u201cSingle-chip CMOS optical microspectrometer\u201d, J. H. Correia et al., Sensors and Actuators, 82 (2000) 1.91-197.\nThese etalon fabrication of fixed cavity filters in monolithically integrated circuits include multiple timed etches for providing varying discrete levels of one mirror. However, establishing the distance of the mirror levels by a timed etch is inaccurate for dimensions involving wavelengths on the order of 400 to 1000 nanometers.\nAccordingly, it is desirable to provide a method of fabricating a Fabry-Perot filter in an integrated circuit whose mirror levels may be accurately determined. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description of the invention and the appended claims, taken in conjunction with the accompanying drawings and this background of the invention."}
-{"text": "1. Field of the Invention\nThe disclosed technology generally relates to single-carrier communication systems wherein adaptive feedback equalization is applied.\n2. Description of the Related Technology\nIn outdoor wireless applications the multipath channel exhibits very long impulse responses, which can lead to delay spreads of several tens of microseconds. This can be of particular importance in e.g. cellular or broadcast applications. When transmitting at high rate over such channels, the introduced intersymbol interference (ISI) can span hundreds of symbols and, hence, severely distorts the received signals. To cope with such ISI, very long decision feedback equalizers (DFE) are required at the receiver to properly recover the transmitted signal. The design of such equalizer can be very complex and occupy a large chip area since certain performance requirements have to be met.\nDecision feedback equalizers (DFE), consisting of a feedforward (FF) and a feedback (FB) filter, are well known in the art. They are preferred to linear equalizers because of their effectiveness at reducing the ISI. This stems mainly from their ability to cancel very efficiently the post-cursor portion of the ISI. This is done by optimally computing the feedforward and feedback coefficients, provided that the channel impulse response is known at the receiver. When this is not the case and the transmission extends over a large number of symbols (as e.g. in broadcasting systems), an adaptive DFE provides the means to compute adaptively the filter coefficients without any prior knowledge about the channel. According to the system definition, this adaptation can be carried out in a data-aided (trained) mode or non data-aided (blind) mode.\nAlthough trained or blind DFEs may provide satisfactory performance, even in high delay spread channels, their implementation complexity grows linearly with the channel length expressed in number of symbol instants. When this length reaches several hundreds of symbol instants, the number of operations can easily exceed a thousand complex multiplications per symbol. For a symbol rate of 10 MHz, which is typical of broadband systems, this translates into 10 billion complex multiplications per seconds, which requires a very high power consumption and/or silicon area. To reduce the burden of the DFE implementation, frequency domain (FD) techniques have been proposed, which enable performing so-called fast convolutions (or correlations), thereby significantly reducing the implementation complexity.\nConsidering the use of FD processing for the linear equalizer is quite straightforward because it is a direct application of the fast convolution or correlation. On the contrary, FD processing for the feedback part of a DFE has been rarely and not always convincingly addressed. The FD DFE presented in the papers \u201cBlind Decision Feedback Equalization for Terrestrial Television Receivers\u201d (M. Ghosh, Proc. IEEE, vol. 86, no. 10, pp. 2070-2081, October 1998) and \u201cOverlap and Save Frequency Domain DFE for Throughput Efficient Single Carrier Transmission\u201d (S. Tomasin, IEEE 16th Int'l Symp. Personal, Indoor and Mobile Radio Communications, September 2005) only have the feedforward part in the FD. The main reason is that the FD block processing does not lend itself easily to the feedback process. The FD approach in \u201cFrequency Domain Feedforward Filter Combined DFE Structure in Single Carrier Systems over Time-varying Channels\u201d (B. Liu et al., IEEE Trans. Cons. Electronics, vol. 54, no. 4, November 2008) targets mainly single-carrier (SC) systems with cyclic extension and is also limited to the FF part. In \u201cFrequency-domain and Multirate Adaptive Filtering\u201d (J. Shynk, IEEE Signal Proc. Magazine, vol. 9, no. 1, pp. 14-37, January 1992) both the FF and FB part, including the weight update, are done in the FD, but the FB is less performant than in the TD case because the new decisions obtained after the DFTs are not fed back until the next DFT is computed.\nRegarding notational conventions, normal Latin characters are used for time-domain signals (a) and tilde characters for frequency-domain signals (\u00e3), vectors and matrices are denoted by a single and double under-bar, respectively, (a and A). The superscripts *, T and H denote the complex conjugate, the matrix transpose and complex conjugate transpose, respectively. The Hadamard (i.e. element-wise) product of vectors is denoted by the \u2299 operator. Any matrix in this description represented by a Q followed by a subscripted letter denotes a zero padding matrix. The operator F represents a FFT block operation.\nBefore a decision feedback equalization scheme as known in the art is described more in detail, a system model is introduced. The transmission is considered of a sequence of possibly complex symbols over a complex multipath channel and affected by additive white Gaussian noise. The discrete-time equivalent of the received signal x(t) can be expressed as:\n                              x          \u2061                      (            k            )                          =                                            \u2211                              l                =                0                                            L                -                1                                      \u2062                                          s                \u2061                                  (                                      k                    -                    l                                    )                                            \u2062                              h                \u2061                                  (                  l                  )                                                              +                      n            \u2061                          (              k              )                                                          (        1        )            where s(l) is the sequence of the transmitted symbols and h(k) denotes the channel impulse response of length L. The received signal x(k) is fed to a DFE (see FIG. 1) that satisfies the following relationship\n                              y          \u2061                      (            k            )                          =                                            \u2211                              i                =                0                            A                        \u2062                                                            a                  i                  *                                \u2061                                  (                  k                  )                                            \u2062                              x                \u2061                                  (                                      k                    +                                          d                      a                                        -                    i                                    )                                                              +                                    \u2211                              i                =                1                            C                        \u2062                                                            c                  i                  *                                \u2061                                  (                  k                  )                                            \u2062                              d                \u2061                                  (                                      k                    -                    i                                    )                                                                                        (        2        )                                          d          \u2061                      (            k            )                          =                  f          \u2062                      {                          y              \u2061                              (                k                )                                      }                                              (        3        )            where ai(k) are the A coefficients of the feedforward filter, ci(k) are the C coefficients of the feedback filter, da is the delay of the feedforward filter and f{.} is the demodulation operator (slicer\u2014in FIG. 1 named \u2018decision device\u2019). In many practical systems, the channel is not known at the receiver and it may be difficult to estimate the channel directly. Adaptive algorithms are then used to adapt the weights ai(k) and ci(k) so as to reduce the variance of the estimation error. Although many high performance algorithms exist, the least-mean square (LMS) algorithm is often used in practice for complexity reasons, especially when the filter lengths A and C are large. The LMS adaptation rule is as follows:e(k)=d(k)\u2212y(k)\u2003\u2003(4)ai(k+1)=ai(k)+2\u03bcax(k\u2212i)e(k)\u2003\u2003(5)ci(k+1)=ci(k)+2\u03bccx(k\u2212i)e(k)\u2003\u2003(6)\nThe LMS algorithm used in combination with the proposed implementation of the DFE provides the simplest adaptation algorithm. However, its complexity increases as the channel length becomes larger. This is because the filtering operations in (2) involve convolutions of length A and C and the updates in (5) and (6) involve correlations of length A and C as well. Since the filter length must be approximately equal to or even greater than the channel length, the total complexity of the LMS-DFE is approximately equal to 4L possibly complex multiplications per output symbol, which may be challenging for channel lengths with hundreds of taps. It has been identified that the frequency domain (FD) processing is an interesting alternative to the time domain (TD) convolutions and correlations as in (2), (5) and (6). Indeed, it is well known that for a length Lx sequence x and length Ly sequence yconv(x,y)=IDFT(DFT(x)\u2299DFT(y))\u2003\u2003(7)corr(x,y)=IDFT(DFT(x)*\u2299DFT(y))\u2003\u2003(8)provided that the discrete Fourier transform (DFT) and inverse discrete Fourier transform (IDFT) are taken over a length larger than or equal to Lx+Ly\u22121. Zeros must be padded to x and y before taking the DFT to achieve this. This alternative implementation is attractive for large Lx and/or Ly if the fast direct and inverse Fourier transforms (FFT and IFFT) are used. When one of the sequences is infinite, it must be split into blocks and special techniques, known as overlap-and-save and overlap-and-add, must be used to recombine the blocks after (7) or (8) is applied.\nHence, there is a need for an adaptive feedback equalization scheme of reasonable complexity, which is suitable for applications where relatively long impulse responses are encountered."}
-{"text": "The present invention relates to controlling temperatures of a cooktop and more particularly to limiting the temperature of a glass-ceramic cooktop such that a maximum temperature of the cooktop reaches about a predetermined maximum temperature.\nThe new trend in electronically controlled cooktops and/or ranges, typically, includes a cooktop surface composed of a glass-ceramic material that is positioned above one or more radiant heating elements. The electronically controlled cooktop includes various user controls that are operated by a user to adjust the amount of heat and, ultimately, the temperature desired for cooking. The radiant heating elements can be powered by electricity, natural gas or other sources. Typically, the radiant heating element and the user controls are connected to a controller that controls the amount of heat supplied to the cooktop. The electronically controlled cooktop also includes temperature and other sensor that are connected to the controller to aid in controlling the heat supplied by the radiant heating source. The temperature and other sensors can also be used in conjunction with the controller to detect certain conditions that can arise during operation of the cooktop.\nTo increase the service life of the glass-ceramic cooktop, the temperature of the glass-ceramic cooktop should not exceed a predetermined maximum temperature for extended periods of time. Typically, the temperature of the glass-ceramic cooktop is affected by several factors during operation. The power level selected by the user using a user control interface is an important factor affecting the temperature of the glass-ceramic cooktop. Generally, the user power selection results in a temperature of the glass-ceramic cooktop that is below the predetermined maximum temperature. However, the user power selection along with other dynamics and/or factors involved with the operation of the glass-ceramic cooktop can cause the temperature of the glass-ceramic cooktop to rise above the predetermined maximum temperature. As mentioned above, exposure to a temperature above the predetermined maximum temperature for extended periods of time may reduce the service life of the glass-ceramic cooktop.\nThe other dynamics and/or factors that affect the glass-ceramic cooktop temperature include, such as, for example, the type of cooking utensil used, the thermal conductivity of the cooking utensil, the contents of the cooking utensil and user manipulation of all of these dynamics and/or factors. The type of cooking utensil affects the temperature of the glass-ceramic cooktop because the size and manufacturer of the cooking utensil can affect the thermal conductivity of the cooking utensil. The thermal conductivity of the cooking utensil relates to the ability of the cooking utensil to transfer heat from the radiant heating element to the contents of the cooking utensil. In addition, factors, such as, for example, the material composition of the cooking utensil, the thickness of the cooking utensil, the flatness and/or warping of the bottom of the cooking utensil also affect the thermal conductivity of the cooking utensil. These factors affect the contact interface area between the cooking utensil and the glass-ceramic cooktop, and the dynamics of this interface area are important factors in determining the thermal conductivity of the cooking utensil. In addition, the contents of the cooking utensil or lack thereof can affect the temperature of the glass-ceramic cooktop because, for example, a cooking utensil that has boiled dry can cause the glass-ceramic cooktop temperature to increase. Additionally, the user can alter or manipulate some or all of these factors/dynamics during the cooking process by, for example, adding more contents to the cooking utensil and/or moving or rocking the cooking utensil on the glass-ceramic cooktop surface during the cooking process.\nAs stated above, various factors and/or dynamics can cause the temperature of the glass-ceramic cooktop to exceed the predetermined maximum temperature, and therefore, potentially reduce the service life of the glass-ceramic cooktop. Therefore, there is a desire to control the power supplied to the radiant heating element such that a maximum temperature of the glass-ceramic cooktop reaches about the predetermined maximum temperature and/or does not exceed the predetermined maximum temperature within accepted tolerance limits."}
-{"text": "The invention relates to a method of compacting a non woven sintered metal web. The invention also relates to a non woven sintered metal web obtained by said method. Such a non woven metal web may be used in filters and in radiant surface combustion burners.\nBy non woven metal web is meant a web consisting of metal fibers having a diameter of less than 100 micrometer, preferably less than 60 micrometer. e.g. 25 .mu.m, 15 .mu.m, 2 .mu.m or 1.5 .mu.m. These metal fibres have rough outer surfaces and are randomly distributed in the web and are interlocked in an intermingled relationship substantially solely by means of the rough outer surfaces thereof in a frictional engagement. A non woven metal web may comprise two or more layers wherein each layer comprises metal fibres of about the same diameter and wherein the diameters of the different layers are not necessarily equal to each other. By a non woven metal web is also meant a non woven metal web reinforced by one or more wire meshes. A non woven metal web may be produced by a rando-feeder-webber apparatus which is disclosed in GB 1.190.844.\nThis non woven metal web may be sintered or not, but relating to this invention the sintered form is to be preferred. By sintering a bond is created at the points of crossing of the different fibers.\nIt is well known to use non woven metal webs in filters. The pore size is one of the most important characteristics of a filter. It determines the size of the particles which will be kept by the filter. In order to obtain a specified limited pore size the non woven metal web has been compacted by pressing or rolling up to now. However, these conventionally compacted non woven webs have some limitations. The density, i.e. a measure For the compactness (for correct definition, see below), corresponding to a specified pore size is quite high. This causes relatively higher pressure drops over the filter or necessitates the use of larger filter surfaces or the use of ventilators and pumps with more power.\nIt is also known to use non woven metal webs as panels in radiant surface combustion burners. These burners can have a true radiant surface combustion, i.e. no blue flame patches or no free flame, with a thermal output up to 900 kW m.sup.-2. However, until now this has only been possible with non woven webs having a limited surface of about 40.000 mm.sup.2 or less (and this not on a constant basis), or, if one desires to use surfaces greater than 40.000 mm.sup.-2, with thermal outputs which are limited to 500 kW m.sup.-2 or less. This represents a severe restriction for the application of these burners."}
-{"text": "Automotive experts quite typically recommend that the squeegees of windshield wiper assemblies be replaced at least annually as a consequence of their deterioration due to abrasion, aging, etc. To minimize the expense of replacement of an entire windshield wiper system including arm, pressure applying superstructure, and squeegee assembly, only the latter is replaced. In the usual case, the squeegee assembly includes both the squeegee itself and a flexor or backing which is applied to the existing pressure applying superstructure mounted on the wiper arm.\nOf course, some means must be provided whereby the squeegee assemblies may be easily disassembled from and replaced on the superstructure which also provides for positive retention of the squeegee assembly on the superstructure during use.\nAs a consequence, there have been a large variety of proposals for retainer clips to accomplish this purpose. Representative of such attempts, and constituting the most pertinent prior art known to the applicant are the following: U.S. Pat. No. 3,626,544 issued Dec. 14, 1971 to Lopez; U.S. Pat. No. 3,879,794 issued Apr. 29, 1975 to Roberts; U.S. Pat. No. 3,919,736 issued Nov. 18, 1975 to Bourassa; U.S. Pat No. 4,009,503 issued Mar. 1, 1977 to Sharp; U.S. Pat. No. 4,156,951 issued June 5, 1979, also to Sharp; U.S. Pat. No. 4,388,742 issued June 21, 1983 to Kimber; and U.S. Pat. No. 4,389,746 issued June 28, 1983 to Riester. Also of possible relevance are French patent publication No. 2,443,354 and German Gebrauchsmuster No. E82 20 739.9. Of the foregoing, U.S. Pat. No. 4,388,742 to Kimber may be the most relevant.\nThe retainers of each of the foregoing patent publications have a variety of advantages and accompanying disadvantages which will be readily appreciated by those skilled in the art. In general, it is desirable that the clip be easily handled, both during an assembly and a disassembly process, that it has sufficient strength as to not be broken or permanently distorted during installation or removal of a squeegee assembly, that it have universal applicability in terms of accommodating superstructures of various configurations and the ability to mount squeegee assemblies of varying length, that it provide for positive retention and firm securement of the squeegee assembly to the superstructure and that it provide a high degree of retention in the proper placement during the assembly process.\nIn general, the prior art suffers from the deficiency of being unable to fulfill one or more of the foregoing desirable attributes of a retaining clip. The present invention is directed to providing a retaining clip that fulfills all of the foregoing desirable attributes."}
-{"text": "1. Field of the Invention\nThe present invention relates to a power conversion controller, especially to a fixed-on-time controller for discontinuous mode PFC (Power Factor Correction) power conversion.\n2. Description of the Related Art\nFixed-on-time power conversion is capable of achieving good PFC for discontinuous mode power conversion when the line input is a full-wave rectified voltage. The reason is as follows:\nLet the line input be VIN(t)=|VA sin(\u03c9t)|, the fixed on time=tON, the inductance=L, and the input current at the end of an on period=IINPEAK(tX), then IINPEAK(tx)=(VIN(tX)/L)tON=(|VA sin(\u03c9tX)|/L)tON, wherein tON is much smaller than 2\u03c0/\u03c9. As can be seen in the above equation, the input current will follow the line input at every on period and a good power factor correction is therefore achieved.\nTo get insight into the principle of the fixed-on-time discontinuous mode PFC power conversion, please refer to FIG. 1, which illustrates a block diagram of a discontinuous mode PFC power converter using a prior art fixed-on-time controller. As illustrated in FIG. 1, the discontinuous mode PFC power converter includes an input rectifier circuit 101, a transformer circuit 102, an output rectifier circuit 103, a load 104, a resistor 105, and a fixed-on-time controller 110.\nThe input rectifier circuit 101 is used to generate a full wave rectified voltage VFULL\u2014WAVE according to an AC power VAC.\nThe transformer circuit 102 is used to transfer the power from the full wave rectified voltage VFULL\u2014WAVE to the output rectifier circuit 103, under the control of a driving signal VDRV of the fixed-on-time controller 110.\nThe output rectifier circuit 103 is used to provide a DC voltage VO to the load 104, and the resistor 105 is used to provide a feedback signal VFB for the fixed-on-time controller 110.\nThe fixed-on-time controller 110 includes an error amplifier 111, a constant current source 112, a capacitor 113, a switch 114, a comparator 115, and a fixed-on-time driver circuit 116.\nThe error amplifier 111 has a negative input end coupled to the feedback signal VFB, a positive input end coupled to a reference voltage VREF, and an output end for providing a threshold signal VCOMP.\nThe constant current source 112, the capacitor 113, and the switch 114 are used to generate a saw signal VSW.\nThe comparator 115 has a negative input end coupled to the threshold signal VCOMP, a positive input end coupled to the saw signal VSW, and an output end for providing a turn-off signal VOFF.\nThe fixed-on-time driver circuit 116 is used to provide the driving signal VDRV for the transformer circuit 102 and a reset signal RESET for the switch 114, according to the turn-off signal VOFF from the comparator 115 and a sensing signal VAUX from the transformer circuit 102, wherein the sensing signal VAUX is used to indicate the end of an inductor current discharging period of the transformer circuit 102, and the active time point of the reset signal RESET follows that of the turn-off signal VOFF.\nWhen in operation, the voltage of the driving signal VDRV will arise from a low level to a high level after the sensing signal VAUX becomes active, and fall from a high level to a low level after the turn-off signal VOFF becomes active. Besides, the period the driving signal VDRV remains at a high level\u2014the on time of the transformer circuit 102\u2014will be fixed to a value by the feedback control mechanism of the power conversion to transfer a specific amount of energy per cycle from the AC power VAC to the load 104, to regulate VFB at VREF.\nHowever, as the period the driving signal VDRV remains at a high level\u2014a fixed on time of the transformer circuit 102 corresponding to a load value of the load 104 and a line voltage of the AC power VAC\u2014is ended when the saw signal VSW reaches the threshold signal VCOMP, the threshold signal VCOMP will exhibit a large level shift to change the period of the fixed on time from a short/long value to a long/short value when the load value of the load 104 or the line voltage of the AC power VAC changes drastically.\nPlease refer to FIG. 2, which illustrates the waveforms of major signals in the prior art fixed-on-time controller 110 of FIG. 1 corresponding to a low line and a high line of the AC power VAC respectively. As illustrated in FIG. 2, VDRV, LOW\u2014LINE (the driving signal VDRV generated corresponding to a low line of the AC power VAC) has a fixed on time tON1, and VDRV, HIGH\u2014LINE (the driving signal VDRV generated corresponding to a high line of the AC power VAC) has a fixed on time tON2, wherein tON1 is longer than tON2 such that same power is delivered to the load 104. As the fixed on time is ended when the saw signal VSW\u2014of which the ramping up slope is fixed by the constant current source 112\u2014reaches the threshold signal VCOMP, the fixed on time is then controlled by the level of the threshold signal VCOMP. As such, a higher VCOMP, LOW\u2014LINE (the threshold signal VCOMP corresponding to the low line of the AC power VAC) is generated to allow a higher VSW, LOW\u2014LINE (the saw signal VSAW corresponding to the low line of the AC power VAC) and therefore the longer tON1; and a lower VCOMP, HIGH\u2014LINE (the threshold signal VCOMP corresponding to the high line of the AC power VAC) is generated to allow a lower VSW, HIGH\u2014LINE (the saw signal VSAW corresponding to the high line of the AC power VAC) and therefore the shorter tON2, wherein the VCOMP, LOW\u2014LINE is higher than the VCOMP, HIGH\u2014LINE by \u0394V1.\nThat is, \u0394V1 of the prior art fixed-on-time controller 110 can be a large value when the load value of the load 104 or the line voltage of the AC power VAC changes drastically. However, large \u0394V1 is adverse to the design of the error amplifier 111. To minimize \u0394V1, one solution is to change the capacitance of the capacitor 113 to cope with the amplitude variation of the AC power VAC or the load value variation of the load 104. However, it is bothersome to change the capacitor 113 whenever the AC power VAC or the load 104 is changed, and besides, the solution is adverse to the integration of the capacitor 113 into the fixed-on-time controller 110.\nIn view of these problems, the present invention proposes a novel fixed-on-time controller for discontinuous mode PFC power conversion."}
-{"text": "Many imaging systems, such as those used in surveillance and targeting systems, include a laser or other light beam projector, such as a laser pointer, laser rangefinder, or laser designator for target designation, rangefinding, or other uses. To ensure accurate aim of the laser beam, the laser beam is typically boresighted (e.g., aligned) to an imaging sensor of the imaging system, for example, to the location of a reticle displayed in images captured by the imaging sensor. For laser designation in military weapons systems, boresighting of laser designators to imaging sensors is of significant importance for accurate delivery of ordinance. However, the alignment between the laser beam and the imaging sensor inevitably drifts over time and temperature.\nTo account for the drift in the alignment, some imaging systems are provided with a boresighting module to perform rapid, automatic boresighting periodically or prior to using the laser beam in the field. However, due to the complex optical arrangement, conventional boresighting modules are typically too complex and bulky for compact applications that fit within and/or selectively move in and out of the optical path of imaging system package. Furthermore, for a target object utilized in boresighting, conventional boresighting modules typically enclose a thermal target that burns away or otherwise deteriorates over use and must be replaced regularly. While imaging systems for surveillance and targeting often include multiples imaging sensors with different wavelength responsivity, conventional boresighting modules typically do not allow for simultaneous boresighting of multiple imaging sensors in imaging systems."}
-{"text": "The invention pertains to a device for projection-copying of a mask onto a workpiece, in particular a semiconductor substrate for the production of integrated circuits, in which case the pattern of the mask is projected via a projection lens onto a photosensitive layer of the workpiece after in order to align the mask and workpiece relative to one another alignment patterns of the mask and adjustment areas on the workpiece have been projected onto each other by means of an adjustment light with a bandwidth of at least 3 nm through the projection lens.\nProjection lenses for the lithographic production of integrated circuits are charcterized by a large picture field with a diameter which typically lies between 10 and 30 mn and a large numerical aperture when there is a great resolution capability which is limited with respect to diffraction. Because of the necessity of projecting different pictures one on top of the other in such a way that they precisely cover one another, the distortion in the entire field must not exceed 0.2 .mu.m and the picture field must be completely level, i.e., any convexity must not exceed 0.3 .mu.m.\nAt the present state of the art, lenses which meet such stringent requirements can be corrected up to the diffraction limit only for a very narrow wavelength range, which is to be understood as a bandwidth of a few nm.\nThis narrow correction area is selected in such a way that within it the sensitivity of the photosensitive resist is as high as possible and, on the other hand, an appropriate source of illumination is available. Typical correction ranges are 406 nm.+-.4 nm or 436 nm.+-.4 nm, corresponding to the most intense spectral lines of mercury discharge lamps.\nAs mentioned above, it is vital for the projection of the pattern of the mask to take place not only with good picture quality, but also with complete precision of positioning. The precision with respect to the lateral coordinates (X, Y, .theta.) is necessary in this case in order for the successive patterns to be correctly allocated, but in addition it is also necessary to precisely focus an entire picture plane since the depth of focus of the above-described lenses is very slight.\nThe alignment of the mask and workpiece is preferably carried out through the lens itself, and in this process the adjustment areas for lateral adjustment are defined by marks of the most widely varying structure. Per se, the unaltered reflective surface of the workpiece itself is sufficient for focusing.\nFocusing as such does not necessarily require that the workpiece be observed through the lens. For instance, it is a known process to determine the distance between the lens and the workpiece by means of capacitive sensors with the aid of the flow resistance which the annular gap between the lens and the workpiece presents to a discharging gas, or from the eigen-frequency of an air resonance section between the workpiece and the lens. A prerequisite in this case, however, is that there be a very short distance between the lens and the workpiece, but such a short distance is precisely what is avoided through the use of projection illumination procedures, in contrast to the obsolete contact procedures. When optical focusing is carried out without the use of the projection lens, a tightly packed beam of light is obliquely directed onto the center of the field to be illuminated and the point where the reflected light beam strikes is observed; the location of this point is a measure of the location of the illumination field itself. The light used is frequently laser light (a HeNe laser). A disadvantage in this case is the fact that it is impossible to distinguish whether the change in the position of the light point on the receiver is caused by a change in the angle of incidence or the Z position. In practice, it is assumed that the angle of incidence is constant, i.e., that the reflective wafer surface is always vertical to the optical axis. This condition is never rigorously met; sometimes there are even large deviations from the ideal position. When laser light is used, another problem arises due to diffraction effects at the wafer surface (speckle), in particular if the surface has already been structured (higher manufacturing stages).\nSince, when adjustment is done through the projection lens, the system-related disadvantages of the above-described procedures do not arise, there is great interest in solving the sub-problems which still exist with this type of adjustment. When adjustment is done through the projection lens, problems arise when the adjustment light is not identical to the illumination light for which the projection lens is corrected. Initially we think here of the case where the wavelength range of the adjustment light lies outside of the range of the spectral sensitivity of the photosensitive resist to avoid having the marks on the workpiece be destroyed by the adjustment process. The differences with respect to focal length and magnification which the lens shows depending on the type of light used can be compensated for by bending the adjustment beam through a pair of mirrors or by lengthening the beam by means of intermediate glasses or by shifting the location of the protector for adjustment. Since tests can readily determine the extent to which the focal length and magnification of the lens differ at the corresponding illumination wavelength and adjustment wavelength, overall the consequences of this difference can be easily handled; with respect to the position where the device is optimally aligned with the adjustment light, prior to illumination being carried out a shift is simply made which takes the differing behavior of the lens in the two cases into account.\nSince, with an illumination wavelength which deviates from the adjustment wavelength, the picture defects of the lens can be corrected only when the bandwidth of the adjustment light is relatively small, it is assumed that the adjustment light used should, in principle, be narrowband. The frequenlty-made suggestion that laser light be used for adjustment purposes has not proven out in practice since, due to the coherence of this light, diffraction effects (speckle) occur which distort the measurement results. Generally mercury discharge lamp light, the natural line width of which is approximately 3 nm, is thus used for adjustment.\nSurprisingly enough, it has been found that, when the adjustment is carried out with mercury light through the projection lens in the way described in the introduction, a nonsystematic error arises, the cause of which was found to be the fact that the reflection capacity of the workpiece is dependent on the processing stage of the workpiece and, in addition, that this capacity varies on the surface of the workpiece. Not only does the color of the reflected adjustment light deviate slightly from the color of the input adjustment light, but differences in the color of the reflected adjustment light appear from workpiece to workpiece and from mark to mark on the same workpiece. If it is assumed that, in addition to the spectral line itself, the adjacent area of the radiation background in a total width of, for example, 10 nm is also typically passed by the narrowband interference filter in front of the adjustment light sources, it is still amazing that the differences of 1-2 nm which arise overall in the wavelength of the reflected adjustment light still have an effect on the precision of the adjustment. Due to the heavy frequency dependency of the lens used, on the one hand, and the extreme demands imposed on adjustment precision, on the other, however, this is indeed the case.\nThe invention thus is based on the recognition that it is not sufficient, as was previously the assumption, to take into account the differences which arise, regarding the focal length and magnification of the lens, in the illumination wavelength on the one hand and, on the other, and adjustment wavelength which is assumed to be constant, but rather it is necessary to take into account the effect of the difference in the reflection behavior of the workpieces, which difference cannot be known in advance, and which leads to a change in the spectral composition of the input adjustment light."}
-{"text": "Integrated circuits are available in many different packages, technologies, and sizes. Most integrated circuits are available in plastic packages, which are generally intended for commercial operating environments at a low cost. Commercial operating environments have a specified operating range from 0\u00b0 C. to 70\u00b0 C. Integrated circuits for military applications have historically been packaged in either metal or ceramic hermetic packages, which are able to work reliably in more demanding environments than commercial integrated circuits. Military operating environments have a specified operating range from \u221255\u00b0 C. to 125\u00b0 C. In order to save costs, the military has purchased integrated circuits through COTS (Commercial Off-The-Shelf) programs. However, these components are generally commercial grade components in plastic packages, and not intended for demanding environments requiring the broader temperature range reliability and durability of ceramic and metal hermetically packaged integrated circuits.\nDepending on size and complexity, integrated circuits are available in a wide range of packages. Although many older integrated circuits were packaged using through-hole technology packages, surface mount packages have dominated over the past several decades. Surface mount packages generally have circuit density, cost, and other advantages over through-hole integrated circuits. Examples of through-hole packages include DIP (dual-in-line plastic) and PGA (pin grid array). Examples of surface mount packages include SOIC (small-outline integrated circuit) and PLCC (plastic leaded chip carrier).\nIntegrated circuit packages generally consist of a semiconductor die placed within a package base and bonded to the base with a suitable die attach adhesive. In conventional technology, the die is electrically attached to a lead frame of the package base with discrete bond wires, which connect individual pads of the die with package leads. In most cases, the bond wires are gold, but in other environments can be copper or aluminum. Specialized equipment is required to attach the bond wires to the die pads the lead frame. Once all of the bond wires are attached, the package lid is bonded to the package base and the integrated circuit can be tested."}
-{"text": "Wire bonding is typically applied to make electrical connections between an integrated circuit die or chip and a carrier on which the die is mounted. Bonding wires are attached to bond pads on the chip and bonding leads on the carrier respectively by ultrasonic welding using an ultrasonic transducer which is integrated into a wire bonding apparatus. The ultrasonic transducer is an energy converting-device which converts electrical energy into ultrasonic vibrations and transmits the ultrasonic vibrations to a capillary at a tip end of the transducer to perform wire bonding.\nFIG. 1 illustrates a side view of a conventional wire bonding apparatus 10 with a transducer 12 positioned over a bonding area for wire bonding. The transducer 12 has a long slender rod-shaped horn 13 and a capillary 16 located at a tip end of the horn 13. Wire is threaded through the capillary 16 and extends out of the capillary tip. The capillary 16 is raised or lowered with respect to a bonding area by moving the transducer 12 vertically.\nAn ultrasonic generator 17 is connected to an end of the horn 13 that is opposite to the tip end of the horn 13 where the capillary 16 is mounted. The ultrasonic generator 17, typically comprising a stack of piezoelectric elements, is housed in a transducer holder 14 mounted on a bond arm 18. The transducer holder 14 is supported by the bond arm 18, which is in turn connected to a sliding bar 22 at a pivot 20. Rotational motion of the bond arm 18 about the pivot 20 rotates the bond arm 18 up and down thereby raising or lowering the transducer 12 relative to the bonding area. The sliding bar 22 may further move and position the bond arm 18 and transducer 12 along the x and y axes.\nA carrier 26, on which components such as semiconductor dice and bonding wires are attached, is held in position by a window clamp 24 on a wire bonding platform to perform wire bonding. The wire bonding platform includes a heater block 28 which provides heat to the carrier 26 to facilitate wire bonding conducted on it.\nDuring wire bonding, the ultrasonic generator 17 produces ultrasonic vibrations which are transmitted along the horn 13. The long slender rod shape of the horn 13 is suitable for amplifying the ultrasonic vibrations transmitted to the capillary 16 while suppressing attenuation. In this way, ultrasonic vibrations may be transferred to the capillary 16 efficiently.\nA pressing force is also applied to the capillary 16, and bonding is accomplished by applying the ultrasonic energy transmitted through the horn 13 onto the wire that is subsequently attached to bonding pads, typically on the semiconductor dice. Control of the positioning of the capillary 16 relative to the bonding area is essential to perform wire bonding accurately. This is particularly so when very fine wires and bond pitches are involved.\nAs can be seen in FIG. 1, the transducer 12 is positioned over the heating zone during wire bonding. Heat is transmitted to the immediately vicinity of the heater block 28, which means that the transducer 12 is exposed to heat from the heater block 28 by radiation. The heat may cause the transducer 12 and the transducer holder 14 to expand. This may therefore affect the relative positioning of the capillary 16 and a bonding position. As a result, wire bonding accuracy is affected.\nA thermal shield may be used for reducing heat transmission from the heater block 28 to the transducer 12 and the transducer holder 14. FIG. 2 is a side view of the wire bonding apparatus 10 with a thermal shield 30 affixed to the sliding bar 22 of the wire bonding apparatus. The thermal shield 30 is operative to separate the transducer 12 from the heater block 28, thereby reducing radiation of heat from the heater block 28 directly to the transducer 12. However, the window clamp 24 needs to be lifted to unclamp the carrier 26 every time the carrier 26 is moved or when it is removed from the heater block 28. The presence of the thermal shield 30 limits the vertical distance A that is movable by the window clamp 24. It is desirable to increase the vertical distance that is movable by the window clamp 24 to ensure sufficient clearance of the window clamp from the carrier 26 and the components attached thereon when the carrier 26 needs to be moved.\nFIG. 3 illustrates a front view of the wire bonding apparatus of FIG. 2 showing the thermal shield 30 between the transducer 12 and the heater block 28. While the thermal shield 30 may reduce the heat radiated from the heater block 28 to the transducer 12, it is noted that the transducer 12 is still exposed to ambient heat from the surrounding air above the thermal shield 30. The transducer 12 and transducer holder 14 would still expand such as to change the relative positioning of the capillary 16 and the carrier 26, resulting in decreased wire bonding accuracy.\nFIG. 4 is a side view of the wire bonding apparatus 10 of FIG. 2 illustrating movement of the transducer 12 along the y and z axes. The thermal shield 30 is movable together with the sliding bar 22 along the x and y axes as it is attached to the sliding bar 22. However, when the transducer 12 moves along the z-axis by rotary movement of the bond arm 18 about the pivot 20, the thermal shield 30 cannot move correspondingly in the z axis. Thus, the transducer 12 is likely to receive more heat energy from the surrounding air as it moves further away from the thermal shield 30. For the aforesaid reasons, use of a conventional thermal shield in a wire bonding system does not effectively address the problem of thermal instability encountered by a transducer during wire bonding. It would be desirable to provide better thermal protection to the transducer 12 from the heater block 28."}
-{"text": "The present invention comprises a new and distinct cultivar of Aglaonema, botanically known as Aglaonema hybrid and hereinafter referred to by the variety name as \u2018TWYAG0042\u2019. The new variety originated from an open pollination made in Bogor, West Java, Indonesia between unknown individual plants of Aglaonema (species unknown). The new variety was discovered as a single plant within the progeny of the stated open pollination in a controlled environment in Bogor, West Java, Indonesia.\nThe new variety was created in Bogor, West Java, Indonesia and has been asexually reproduced repeatedly by vegetative cuttings and tissue culture in Apopka, Fla. for two or more generations. The present invention has been found to retain its distinctive characteristics through successive asexual propagations."}