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base_model: stabilityai/stable-diffusion-xl-base-1.0
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license: cc-by-nc-nd-4.0
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---
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<Gallery />
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## Trigger words
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base_model: stabilityai/stable-diffusion-xl-base-1.0
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license: cc-by-nc-nd-4.0
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---
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# ⚡ FlashDiffusion: FlashSD ⚡
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Flash Diffusion is a diffusion distillation method proposed in [ADD ARXIV]() *by Clément Chadebec, Onur Tasar and Benjamin Aubin.*
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This model is a **26.4M** LoRA distilled version of SD1.5 model that is able to generate images in **2-4 steps**. The main purpose of this model is to reproduce the main results of the paper.
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<p align="center">
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<img style="width:700px;" src="images/hf_grid.png">
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</p>
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# How to use?
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The model can be used using the `StableDiffusionPipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **2-4 steps**.
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```python
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from diffusers import DiffusionPipeline, LCMScheduler
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adapter_id = "jasperai/flash-sd"
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pipe = DiffusionPipeline.from_pretrained(
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"stabilityai/stable-diffusion-xl-base-1.0",
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use_safetensors=True,
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)
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pipe.scheduler = LCMScheduler.from_pretrained(
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"stabilityai/stable-diffusion-xl-base-1.0",
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subfolder="scheduler",
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timestep_spacing="trailing",
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)
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pipe.to("cuda")
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# Fuse and load LoRA weights
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pipe.load_lora_weights(adapter_id)
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pipe.fuse_lora()
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prompt = "A raccoon reading a book in a lush forest."
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image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0]
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```
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<p align="center">
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<img style="width:400px;" src="images/raccoon.png">
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</p>
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# Training Details
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## License
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This model is released under the the Creative Commons BY-NC license.
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