Update README.md

README.md

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 # ⚡ Flash Diffusion: FlashSDXL ⚡
 
 
-Flash Diffusion is a diffusion distillation method proposed in [ADD ARXIV]() *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin.*
+Flash Diffusion is a diffusion distillation method proposed in [Flash Diffusion: Accelerating Any Conditional
+Diffusion Model for Few Steps Image Generation](http://arxiv.org/abs/2406.02347) *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin.*
 This model is a **108M** LoRA distilled version of [SDXL](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0) model that is able to generate images in **4 steps**. The main purpose of this model is to reproduce the main results of the paper.
 See our [live demo](https://huggingface.co/spaces/jasperai/FlashPixart).
 
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 # Training Details
-The model was trained for 20k iterations on 4 H100 GPUs (representing approximately a total of 176 GPU hours of training). Please refer to the [paper]() for further parameters details. 
+The model was trained for 20k iterations on 4 H100 GPUs (representing approximately a total of 176 GPU hours of training). Please refer to the [paper](http://arxiv.org/abs/2406.02347) for further parameters details. 
 
 **Metrics on COCO 2014 validation (Table 3)**
   - FID-10k: 21.62 (4 NFE)