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Update app.py
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app.py
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@@ -3,46 +3,31 @@ import torch
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import numpy as np
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import modin.pandas as pd
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from PIL import Image
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from diffusers import DiffusionPipeline #, StableDiffusion3Pipeline
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from huggingface_hub import hf_hub_download
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device = 'cuda' if torch.cuda.is_available() else 'cpu'
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torch.cuda.max_memory_allocated(device=device)
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torch.cuda.empty_cache()
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generator = np.random.seed(0) if seed == 0 else torch.manual_seed(seed)
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pipe = DiffusionPipeline.from_pretrained("circulus/canvers-real-v3.9.1", torch_dtype=torch.float16, safety_checker=None) if torch.cuda.is_available() else DiffusionPipeline.from_pretrained("circulus/canvers-real-v3.9.1")
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pipe.enable_xformers_memory_efficient_attention()
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pipe = pipe.to(device)
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torch.cuda.empty_cache()
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image = pipe(Prompt, negative_prompt=negative_prompt, height=height, width=width, num_inference_steps=steps, guidance_scale=scale).images[0]
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torch.cuda.empty_cache()
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return image
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if Model == "Animagine XL 4":
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animagine = DiffusionPipeline.from_pretrained("cagliostrolab/animagine-xl-4.0", torch_dtype=torch.float16, safety_checker=None) if torch.cuda.is_available() else DiffusionPipeline.from_pretrained("cagliostrolab/animagine-xl-4.0")
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animagine.enable_xformers_memory_efficient_attention()
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animagine = animagine.to(device)
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torch.cuda.empty_cache()
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image = animagine(Prompt, negative_prompt=negative_prompt, height=height, width=width, num_inference_steps=steps, guidance_scale=scale).images[0]
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torch.cuda.empty_cache()
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return image
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return image
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gr.Interface(fn=genie, inputs=[gr.Radio(['PhotoReal', 'Animagine XL 4',], value='PhotoReal', label='Choose Model'),
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gr.Textbox(label='What you want the AI to generate. 77 Token Limit.'),
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gr.Textbox(label='What you Do Not want the AI to generate. 77 Token Limit'),
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gr.Slider(512, 1024, 768, step=128, label='Height'),
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gr.Slider(512, 1024, 768, step=128, label='Width'),
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gr.Slider(3, maximum=12, value=5, step=.25, label='Guidance Scale', info="5-7 for PhotoReal and 7-10 for Animagine"),
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gr.Slider(25, maximum=50, value=25, step=25, label='Number of Iterations'),
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gr.Slider(minimum=0, step=1, maximum=9999999999999999, randomize=True, label='Seed: 0 is Random'),
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],
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outputs=gr.Image(label='Generated Image'),
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import numpy as np
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import modin.pandas as pd
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from PIL import Image
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from diffusers import StableDiffusion3Pipeline #DiffusionPipeline #, StableDiffusion3Pipeline
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from huggingface_hub import hf_hub_download
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device = 'cuda' if torch.cuda.is_available() else 'cpu'
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torch.cuda.max_memory_allocated(device=device)
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torch.cuda.empty_cache()
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pipe = StableDiffusion3Pipeline.from_pretrained("stabilityai/stable-diffusion-3.5-large-turbo", torch_dtype=torch.bfloat16)
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pipe = pipe.to(device)
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def genie (Prompt, height, width, seed):
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generator = np.random.seed(0) if seed == 0 else torch.manual_seed(seed)
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image = pipe(Prompt, num_inference_steps=4, height=height, width=width, guidance_scale=0.0,).images[0]
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return image
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gr.Interface(fn=genie, inputs=[gr.Radio(['PhotoReal', 'Animagine XL 4',], value='PhotoReal', label='Choose Model'),
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gr.Textbox(label='What you want the AI to generate. 77 Token Limit.'),
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#gr.Textbox(label='What you Do Not want the AI to generate. 77 Token Limit'),
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gr.Slider(512, 1024, 768, step=128, label='Height'),
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gr.Slider(512, 1024, 768, step=128, label='Width'),
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#gr.Slider(3, maximum=12, value=5, step=.25, label='Guidance Scale', info="5-7 for PhotoReal and 7-10 for Animagine"),
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#gr.Slider(25, maximum=50, value=25, step=25, label='Number of Iterations'),
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gr.Slider(minimum=0, step=1, maximum=9999999999999999, randomize=True, label='Seed: 0 is Random'),
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],
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outputs=gr.Image(label='Generated Image'),
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