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-# Stable Diffusion web UI
-A browser interface based on Gradio library for Stable Diffusion.
-
-![](txt2img_Screenshot.png)
-
-Check the [custom scripts](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Custom-Scripts) wiki page for extra scripts developed by users.
-
-## Features
-[Detailed feature showcase with images](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Features):
-- Original txt2img and img2img modes
-- One click install and run script (but you still must install python and git)
-- Outpainting
-- Inpainting
-- Color Sketch
-- Prompt Matrix
-- Stable Diffusion Upscale
-- Attention, specify parts of text that the model should pay more attention to
-    - a man in a ((tuxedo)) - will pay more attention to tuxedo
-    - a man in a (tuxedo:1.21) - alternative syntax
-    - select text and press ctrl+up or ctrl+down to automatically adjust attention to selected text (code contributed by anonymous user)
-- Loopback, run img2img processing multiple times
-- X/Y plot, a way to draw a 2 dimensional plot of images with different parameters
-- Textual Inversion
-    - have as many embeddings as you want and use any names you like for them
-    - use multiple embeddings with different numbers of vectors per token
-    - works with half precision floating point numbers
-    - train embeddings on 8GB (also reports of 6GB working)
-- Extras tab with:
-    - GFPGAN, neural network that fixes faces
-    - CodeFormer, face restoration tool as an alternative to GFPGAN
-    - RealESRGAN, neural network upscaler
-    - ESRGAN, neural network upscaler with a lot of third party models
-    - SwinIR and Swin2SR([see here](https://github.com/AUTOMATIC1111/stable-diffusion-webui/pull/2092)), neural network upscalers
-    - LDSR, Latent diffusion super resolution upscaling
-- Resizing aspect ratio options
-- Sampling method selection
-    - Adjust sampler eta values (noise multiplier)
-    - More advanced noise setting options
-- Interrupt processing at any time
-- 4GB video card support (also reports of 2GB working)
-- Correct seeds for batches
-- Live prompt token length validation
-- Generation parameters
-     - parameters you used to generate images are saved with that image
-     - in PNG chunks for PNG, in EXIF for JPEG
-     - can drag the image to PNG info tab to restore generation parameters and automatically copy them into UI
-     - can be disabled in settings
-     - drag and drop an image/text-parameters to promptbox
-- Read Generation Parameters Button, loads parameters in promptbox to UI
-- Settings page
-- Running arbitrary python code from UI (must run with --allow-code to enable)
-- Mouseover hints for most UI elements
-- Possible to change defaults/mix/max/step values for UI elements via text config
-- Random artist button
-- Tiling support, a checkbox to create images that can be tiled like textures
-- Progress bar and live image generation preview
-- Negative prompt, an extra text field that allows you to list what you don't want to see in generated image
-- Styles, a way to save part of prompt and easily apply them via dropdown later
-- Variations, a way to generate same image but with tiny differences
-- Seed resizing, a way to generate same image but at slightly different resolution
-- CLIP interrogator, a button that tries to guess prompt from an image
-- Prompt Editing, a way to change prompt mid-generation, say to start making a watermelon and switch to anime girl midway
-- Batch Processing, process a group of files using img2img
-- Img2img Alternative, reverse Euler method of cross attention control
-- Highres Fix, a convenience option to produce high resolution pictures in one click without usual distortions
-- Reloading checkpoints on the fly
-- Checkpoint Merger, a tab that allows you to merge up to 3 checkpoints into one
-- [Custom scripts](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Custom-Scripts) with many extensions from community
-- [Composable-Diffusion](https://energy-based-model.github.io/Compositional-Visual-Generation-with-Composable-Diffusion-Models/), a way to use multiple prompts at once
-     - separate prompts using uppercase `AND`
-     - also supports weights for prompts: `a cat :1.2 AND a dog AND a penguin :2.2`
-- No token limit for prompts (original stable diffusion lets you use up to 75 tokens)
-- DeepDanbooru integration, creates danbooru style tags for anime prompts (add --deepdanbooru to commandline args)
-- [xformers](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Xformers), major speed increase for select cards: (add --xformers to commandline args)
-- History tab: view, direct and delete images conveniently within the UI
-- Generate forever option
-- Training tab
-     - hypernetworks and embeddings options
-     - Preprocessing images: cropping, mirroring, autotagging using BLIP or deepdanbooru (for anime)
-- Clip skip
-- Use Hypernetworks
-- Use VAEs
-- Estimated completion time in progress bar
-- API
-- Support for dedicated [inpainting model](https://github.com/runwayml/stable-diffusion#inpainting-with-stable-diffusion) by RunwayML. 
-- Aesthetic Gradients, a way to generate images with a specific aesthetic by using clip images embds (implementation of [https://github.com/vicgalle/stable-diffusion-aesthetic-gradients](https://github.com/vicgalle/stable-diffusion-aesthetic-gradients))
-
-
-## Installation and Running
-Make sure the required [dependencies](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Dependencies) are met and follow the instructions available for both [NVidia](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Install-and-Run-on-NVidia-GPUs) (recommended) and [AMD](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Install-and-Run-on-AMD-GPUs) GPUs.
-
-Alternatively, use online services (like Google Colab):
-
-- [List of Online Services](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Online-Services)
-
-### Automatic Installation on Windows
-1. Install [Python 3.10.6](https://www.python.org/downloads/windows/), checking "Add Python to PATH"
-2. Install [git](https://git-scm.com/download/win).
-3. Download the stable-diffusion-webui repository, for example by running `git clone https://github.com/AUTOMATIC1111/stable-diffusion-webui.git`.
-4. Place `model.ckpt` in the `models` directory (see [dependencies](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Dependencies) for where to get it).
-5. _*(Optional)*_ Place `GFPGANv1.4.pth` in the base directory, alongside `webui.py` (see [dependencies](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Dependencies) for where to get it).
-6. Run `webui-user.bat` from Windows Explorer as normal, non-administrator, user.
-
-### Automatic Installation on Linux
-1. Install the dependencies:
-```bash
-# Debian-based:
-sudo apt install wget git python3 python3-venv
-# Red Hat-based:
-sudo dnf install wget git python3
-# Arch-based:
-sudo pacman -S wget git python3
-```
-2. To install in `/home/$(whoami)/stable-diffusion-webui/`, run:
-```bash
-bash <(wget -qO- https://raw.githubusercontent.com/AUTOMATIC1111/stable-diffusion-webui/master/webui.sh)
-```
-
-### Installation on Apple Silicon
-
-Find the instructions [here](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Installation-on-Apple-Silicon).
-
-## Contributing
-Here's how to add code to this repo: [Contributing](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki/Contributing)
-
-## Documentation
-The documentation was moved from this README over to the project's [wiki](https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki).
-
-## Credits
-- Stable Diffusion - https://github.com/CompVis/stable-diffusion, https://github.com/CompVis/taming-transformers
-- k-diffusion - https://github.com/crowsonkb/k-diffusion.git
-- GFPGAN - https://github.com/TencentARC/GFPGAN.git
-- CodeFormer - https://github.com/sczhou/CodeFormer
-- ESRGAN - https://github.com/xinntao/ESRGAN
-- SwinIR - https://github.com/JingyunLiang/SwinIR
-- Swin2SR - https://github.com/mv-lab/swin2sr
-- LDSR - https://github.com/Hafiidz/latent-diffusion
-- Ideas for optimizations - https://github.com/basujindal/stable-diffusion
-- Doggettx - Cross Attention layer optimization - https://github.com/Doggettx/stable-diffusion, original idea for prompt editing.
-- InvokeAI, lstein - Cross Attention layer optimization - https://github.com/invoke-ai/InvokeAI (originally http://github.com/lstein/stable-diffusion)
-- Rinon Gal - Textual Inversion - https://github.com/rinongal/textual_inversion (we're not using his code, but we are using his ideas).
-- Idea for SD upscale - https://github.com/jquesnelle/txt2imghd
-- Noise generation for outpainting mk2 - https://github.com/parlance-zz/g-diffuser-bot
-- CLIP interrogator idea and borrowing some code - https://github.com/pharmapsychotic/clip-interrogator
-- Idea for Composable Diffusion - https://github.com/energy-based-model/Compositional-Visual-Generation-with-Composable-Diffusion-Models-PyTorch
-- xformers - https://github.com/facebookresearch/xformers
-- DeepDanbooru - interrogator for anime diffusers https://github.com/KichangKim/DeepDanbooru
-- Initial Gradio script - posted on 4chan by an Anonymous user. Thank you Anonymous user.
-- (You)