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| from diffusers import StableDiffusionPipeline | |
| import torch | |
| from PIL import Image | |
| # Check if CUDA is available | |
| device = "cuda" if torch.cuda.is_available() else "cpu" | |
| # Load the pre-trained Stable Diffusion 3.5 model | |
| pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-3.5-large-turbo") | |
| pipe.to(device) | |
| # If you have LoRA weights, load them (replace with your actual LoRA weight path) | |
| pipe.load_lora_weights("prithivMLmods/SD3.5-Turbo-Realism-2.0-LoRA") | |
| # Text prompt for image generation | |
| prompt = "Turbo Realism, High-resolution photograph, woman, UHD, photorealistic, shot on a Sony A7III --chaos 20 --ar 1:2 --style raw --stylize 250" | |
| # Generate the image | |
| with torch.no_grad(): # Disable gradient calculations for inference | |
| image = pipe(prompt).images[0] | |
| # Show the generated image | |
| image.show() | |
| # Optionally, save the image to disk | |
| image.save("generated_image.png") | |