Text-to-Image
Diffusers
TensorBoard
stable-diffusion
diffusers-training
lora
template:sd-lora
stable-diffusion-xl
stable-diffusion-xl-diffusers
import torch
from diffusers import DiffusionPipeline
# switch to "mps" for apple devices
pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", dtype=torch.bfloat16, device_map="cuda")
pipe.load_lora_weights("Helith/sdxl-base-1.0-kavya")
prompt = "portrait of [v]kavya"
image = pipe(prompt).images[0]SDXL LoRA DreamBooth - Helith/sdxl-base-1.0-kavya
- Prompt
- portrait of [v]kavya
- Prompt
- portrait of [v]kavya
- Prompt
- portrait of [v]kavya
- Prompt
- portrait of [v]kavya
Model description
These are Helith/sdxl-base-1.0-kavya LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.
The weights were trained using DreamBooth.
LoRA for the text encoder was enabled: True.
Special VAE used for training: madebyollin/sdxl-vae-fp16-fix.
Trigger words
You should use [v]kavya to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
Intended uses & limitations
How to use
# TODO: add an example code snippet for running this diffusion pipeline
Limitations and bias
[TODO: provide examples of latent issues and potential remediations]
Training details
[TODO: describe the data used to train the model]
- Downloads last month
- 11
Model tree for Helith/sdxl-base-1.0-kavya
Base model
stabilityai/stable-diffusion-xl-base-1.0