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πŸŒƒ SDXL-Lightning-8-Steps (HyHorX Edition)

Developed by: HyHorX AI Lab
Model Type: Faster Diffusion-based text-to-image generation model
License: CreativeML Open RAIL-M

This model is a high-speed merge of Stable Diffusion XL 1.0 Base and ByteDance's SDXL-Lightning 8-step UNet. It is optimized for ultra-fast generation while maintaining high-quality architectural and cinematic details.

πŸš€ Usage Details

To get the best results (like the Cyberpunk city preview), use the following settings:

  • Inference Steps: 8 steps (Strictly)
  • Guidance Scale (CFG): 0 (or 1.0)
  • Sampler/Scheduler: Euler Ancestral (Euler A)
  • Timestep Spacing: Trailing

πŸ’» Example Code (Diffusers)

import torch
from diffusers import StableDiffusionXLPipeline

pipe = StableDiffusionXLPipeline.from_pretrained(
    "HyHorX/SDXL-Lighting-8-steps", 
    torch_dtype=torch.float16, 
    use_safetensors=True
).to("cuda")

prompt = "A majestic cinematic view of a futuristic city with neon lights, high detail, 8k"

image = pipe(
    prompt=prompt, 
    num_inference_steps=8, 
    guidance_scale=0
).images[0]

image.save("result.png")

πŸ› οΈ Optimization Lab
This model was merged and exported in fp16 format to be lightweight (~6.95 GB), making it accessible for consumer GPUs with lower VRAM (like GTX 1060 6GB/8GB).

Created with ❀️ by HyHorX AI Lab during a deep-learning optimization session.

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