Text-to-Image
Diffusers
Trained with AutoTrain
stable-diffusion-xl
stable-diffusion-xl-diffusers
lora
template:sd-lora
import torch
from diffusers import DiffusionPipeline
# switch to "mps" for apple devices
pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", dtype=torch.bfloat16, device_map="cuda")
pipe.load_lora_weights("LinAnnJose/Oldmancharactergeneration")
prompt = "photo of charctr_omn with a plain green background"
image = pipe(prompt).images[0]AutoTrain SDXL LoRA DreamBooth - LinAnnJose/Oldmancharactergeneration
- Prompt
- photo of charctr_omn with a plain green background
- Prompt
- photo of charctr_omn with a plain green background
- Prompt
- photo of charctr_omn with a plain green background
- Prompt
- photo of charctr_omn with a plain green background
Model description
These are LinAnnJose/Oldmancharactergeneration LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.
The weights were trained using DreamBooth.
LoRA for the text encoder was enabled: False.
Special VAE used for training: None.
Trigger words
You should use photo of charctr_omn with a plain green background to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
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Model tree for LinAnnJose/Oldmancharactergeneration
Base model
stabilityai/stable-diffusion-xl-base-1.0