|
|
| --- |
| license: mit |
| base_model: runwayml/stable-diffusion-v1-5 |
| tags: |
| - stable-diffusion |
| - text-to-image |
| - diffusers |
| - lora |
| - style |
| - art |
| --- |
| |
| # LoRA for Style: Impressionism |
|
|
| A LoRA fine-tuned to generate images in the **Impressionism** style. Use the trigger phrase `in impressionism style` in your prompts. |
|
|
| This LoRA was trained on the [WikiArt dataset](https://www.kaggle.com/datasets/steubk/wikiart) for 1000 steps. |
|
|
| ## How to Use |
|
|
| This LoRA is designed to be used with the `diffusers` library. |
|
|
| ```python |
| import torch |
| from diffusers import StableDiffusionPipeline |
| from peft import PeftModel |
| |
| # --- 1. Load the base model --- |
| base_model_id = "runwayml/stable-diffusion-v1-5" |
| pipe = StableDiffusionPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16).to("cuda") |
| |
| # --- 2. Load the LoRA adapter --- |
| repo_id = "aanchal77/Final-One" |
| lora_path_in_repo = "styles/Impressionism" |
| pipe.unet = PeftModel.from_pretrained(pipe.unet, repo_id, subfolder=lora_path_in_repo) |
| |
| # --- 3. Generate an image --- |
| prompt = f"A futuristic cityscape in impressionism style" |
| image = pipe(prompt, num_inference_steps=20, guidance_scale=7.5).images[0] |
| image.save("generated_image.png") |
| |