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Windows is an event-based , or messaging, operating system. This means that Windows applications communicate by sending messages to controls that have Windows handles, and controls that have Windows handles can forward messages to constituent controls. A few years ago programmers had to handle raw messages. Later, in an evolutionary step, raw messages were converted to method invocations. The messages are commonly referred to as events , and the methods that respond are dubbed event handlers . An event handler is quite literally a method designated to respond to message events. One of the most common tasks you will perform is to associate event handlers with control events and then write code for those handlers. This is very similar to how to we created event handlers in VB6. We'll begin here to find some familiar ground and then proceed by exploring what has been revised. Using the Form Designer The easiest way to generate an event handler is to double-click on a control. The Form Designer will generate the event handler designated as the default event. For example, if you double-click on a form, the Load event handler will be generated, and focus will be switched to the Code Designer view. So far this is exactly like VB6. However, there are differences. The signature of the Load event handler looks different, and the event handler is wired to the control differently. Listing 3.1 shows a VB .NET form that contains just the Load event handler. Listing 3.1 Code for a VB .NET Form with the Generated Load Event Handler 1: Public Class Form1 2: Inherits System.Windows.Forms.Form 3: 4: #Region " Windows Form Designer generated code " 5: 6: Public Sub New() 7: MyBase.New() 8: 9: 'This call is required by the Windows Form Designer. 10: InitializeComponent() 11: 12: 'Add any initialization after the InitializeComponent() call 13: 14: End Sub 15: 16: 'Form overrides dispose to clean up the component list. 17: Protected Overloads Overrides Sub Dispose(ByVal disposing As Boolean) 18: If disposing Then 19: If Not (components Is Nothing) Then 20: components.Dispose() 21: End If 22: End If 23: MyBase.Dispose(disposing) 24: End Sub 25: 26: 'Required by the Windows Form Designer 27: Private components As System.ComponentModel.IContainer 28: 29: 'NOTE: The following procedure is required by the Windows Form Designer 30: 'It can be modified using the Windows Form Designer. 31: 'Do not modify it using the code editor. 32: <System.Diagnostics.DebuggerStepThrough()> _ 33: Private Sub InitializeComponent() 34: ' 35: 'Form1 36: ' 37: Me.AutoScaleBaseSize = New System.Drawing.Size(6, 15) 38: Me.ClientSize = New System.Drawing.Size(292, 260) 39: Me.Name = "Form1" 40: Me.Text = "Form1" 41: 42: End Sub 43: 44: #End Region 45: 46: Private Sub Form1_Load(ByVal sender As System.Object, _ 47: ByVal e As System.EventArgs) Handles MyBase.Load 48: 49: End Sub 50: End Class At first glance this seems like a lot of code. However, if you look at the code in a VB6 .frm file (see Listing 3.2), you'll see that the VB .NET code is not that much longer. Listing 3.2 Code for a VB6 Form with the Generated Load Event Handler 1: VERSION 5.00 2: Begin VB.Form Form1 3: Caption = "Form1" 4: ClientHeight = 2400 5: ClientLeft = 48 6: ClientTop = 432 7: ClientWidth = 3744 8: LinkTopic = "Form1" 9: ScaleHeight = 2400 10: ScaleWidth = 3744 11: StartUpPosition = 3 'Windows Default 12: End 13: Attribute VB_Name = "Form1" 14: Attribute VB_GlobalNameSpace = False 15: Attribute VB_Creatable = False 16: Attribute VB_PredeclaredId = True 17: Attribute VB_Exposed = False 18: Option Explicit 19: 20: Private Sub Form_Load() 21: 22: End Sub I digress. Back to the code in Listing 3.1, which explicitly defines a class by using the class header (line 1) and the End Class statement (line 50). A constructor appears in lines 6 through 14 and a destructor in lines 17 through 24, and lines 37 through 40 contain form initialization code. Note that all the code from lines 4 to 44 is managed by the Form Designer. Notice the signature of the Load event handler in lines 46 and 47 of Listing 3.1. Load has two parameters and a Handles clause at the end of the procedure header. The Handles clause indicates that this method ” Handles MyBase.Load ”handles the base class's Load event. (Refer to the upcoming subsection on the Handles clause for more information on this keyword and to the subsection on EventHandler for more information on the arguments passed to the Load event.) Using the Form Designer is the easiest way to generate event handlers. You already know how to click on controls and add code; we won't elaborate on this point any longer. Using the Code Editor The next step familiar to VB6 developers is to generate events in the Code Editor (or Code Designer). The Code Editor has a list of class names at the top left of the editor and a list of method names at the top right. Select a class and a method, and .NET will generate any event handler available for you. It is important to keep in mind that inheritance is an integral part of VB .NET. If you look at the list of method names for the Form1 control (see Figure 3.1), you are likely to be disappointed ”the list is very short. Figure 3.1. A list showing method names available for the selected class ( Form1 ). There is a reason the list is short. We are looking at the events for Form1 . Form1 is derived from Form , in which all the events are defined. If you select Base Class Events in the list of class names under Form1 , you will quickly see that the list of method names has increased substantially. Select Base Class Events from the class name list and the Click method from the method name list to generate a Click event handler. You will quickly see that understanding the impact inheritance has on event name filtering is all you need to know to generate events in the Code Editor. The Handles Clause Visual Basic .NET no longer hides the mechanism that wires events to event handlers when those events are generated by the Form Designer or the Code Editor. The mechanism for this is the Handles clause. Because the Handles clause is a prominent part of the code, it should be apparent that you can write the code manually to handle events at design time. For instance, if we want the Form Load and Click handlers to respond in the same way, we can simply add a MyBase.Click event to the comma-delimited list of events a particular method handles. Listing 3.3 shows the revision. Listing 3.3 A Single Event Handler for Handling Multiple Events Private Sub Form1_Load(ByVal sender As System.Object, _ ByVal e As System.EventArgs) Handles MyBase.Load, MyBase.Click MsgBox("Called on Load and Click") End Sub The Handles clause is now understood to mean that Form1_Load will handle both the Form.Load and Form.Click events. This will work as long as the signatures of the events are identical. That is, one method can respond for any or all methods whose events have identical signatures. (I have not tested the limits of the number of events a single method can respond to or how long the Handles clause can be, but based on the implementation, it should be a very large number of events.) The EventHandler Class The signature of the Load and Click events (and many basic events in .NET) is a method that takes an object and a System.EventArgs argument. (This is actually the signature of a delegate class named EventHandler . We'll come back to the subject of delegates in the What Are Delegates? section later in this chapter.) For now let's explore the arguments of the signature of this basic event handler. The sender Argument You have read by now that classes in Visual Basic .NET are derived from the Object class. Similarly, sender is defined as an object. This means that literally any type can be passed to satisfy the sender argument. What is usually passed is the originating object. The result is that you can use runtime type checking to determine the class of the object that invoked a particular event. Of course, you could write code that depends on the specific object to which you assigned the event handler, but this is less robust, and if you change the name of the object, the code will be wrong. Referring to Listing 3.3, we could revise the code to refer to the object we know is associated with the Form1_Load handler. Private Sub Form1_Load(ByVal sender As System.Object, _ ByVal e As System.EventArgs) Handles MyBase.Load, MyBase.Click MsgBox(Me.Name & " called on Load and Click") End Sub The revised code relies on the reference to self, Me . Specifically, the code relies on the method being associated with the form. How can we implement the code to be less dependent on the specific type of the object? Or how can we write the code to respond differently based on the actual invoker of the event? The answer is that we have to write the code to query the type of the sender object. Private Sub Form1_Load(ByVal sender As System.Object, _ ByVal e As System.EventArgs) Handles MyBase.Load, MyBase.Click If (TypeOf sender Is Form1) Then MsgBox(CType(sender, Form1).Name & " called on Load and Click") End If End Sub The revised code is more robust. The If statement checks to see if the type of the sender object is an instance of Form1 . If it is, sender is dynamically cast (cast at runtime) to its actual type. Casting the sender object supports accessing the members of that type. (For instance, object does not have a name property. However, if sender is an instance of Form1 and we cast it to Form1 , we can access the Form.Name property.) By using dynamic type checking we take specific action based on the real type of the object that raised the event. The EventArgs Argument The EventArgs argument is a stub. In some instances events will want to pass additional information. For example, a mouse event may want to pass the position of the mouse. Keyboard events may want to pass the state of the keyboard. Paint events may want to pass the device context. Other kinds of events actually use the EventArgs class, or derived classes, to pass additional information to the event handler. The Paint event handler passes a PaintEventArgs object that has a Graphics property representing the device context, or canvas, of the invoking object. You can use this Graphics object to perform custom painting, as we did in Chapter 2. Handling Multiple Events with One Handler You now know that you can manually associate similarly signatured events to a single event handler. You can accomplish this task by adding the name of the control and the name of the event in a Handles clause. You may wonder why you would do this. The answer is that you have more than one metaphor that performs the same operation. (A control is a metaphor for invoking an operation.) For example, think of a form that has a menu that closes the form as well as a button that closes the form. Both the menu and the button are metaphors for closing the form. Thus, in addition to performing the same operation, both metaphors should share the same code. The short fragment that follows demonstrates this technique. Private Sub ButtonClose_Click(ByVal sender As System.Object, _ ByVal e As System.EventArgs) _ Handles ButtonClose.Click, MenuItemExit.Click Close() End Sub The event handler was generated by double-clicking the button. From the Handles clause it is clear that this event handler handles the ButtonClose.Click event. I manually typed MenuItemExit.Click to add code for the menu item. The implication is that the user can perform the same operation ”closing the form ”by clicking the button or selecting the menu item. A reasonable person might argue that the savings are minimal. In this example, yes, but when coding we have the opportunity to build habits that work efficiently and effectively in all circumstances. By sharing events that perform identical operations, we converge the code to one function. By calling a well-named function, the code does not need a comment, and we have to change the code in only one place if revisions are needed. To reiterate, we get speed, extensibility, and reliability by building solid general habits and employing them dogmatically. The preceding fragment demonstrates the following habits: Convergent, well-named code helps programmers code at a much faster pace and yields a high-quality result. The objective is to make this style of programming second nature by practicing good habits. Implementing Multiple Respondents You know that you can make one method respond to multiple events. You can also make multiple event handlers respond to a single event. When an event invokes more than one method, this is referred to as multicasting . By implementing multiple event handlers it is possible to layer behavior. For example, a producer can implement a custom Paint event handler for a user control, and a consumer of that control can implement a custom event handler too without overwriting the event handler associated with the Paint event by the producer. (We will return to this subject and provide examples in an upcoming section, What Are Delegates?)
Do you have a student who is finding a particular math skill challenging? Could you use some guidance for teaching a certain concept? The Math Center provides easy access to many of the math-related resources available on SchoolhouseTeachers.com. You’ll find videos, games, eBooks, and other instructional materials about fractions, geometry, addition, decimals, division, and more. This 36-week Algebra 1 course is designed for students who have completed pre-algebra and are ready to advance to the next level of math. This course covers a pre-algebra review, algebraic expressions, solving linear equations, graphing, inequalities, systems of equations, polynomials, factoring, and exponential expressions. It concludes with a final test. Algebra 2 is a full-year course designed for students who have successfully completed Algebra 1. It includes textbook reading, practice problems, homework, and unit tests. Key vocabulary is highlighted, and study tips are included. An answer key is also included. All About Shapes is a fun way to teach and reinforce basic shapes with your preschooler through games, foods, books, crafts, and more! Learn about circles, ovals, squares, rectangles, triangles, stars, pentagons, octagons, crescents, hearts, and diamonds/rhombi. Building a Foundation with Kindergarten Math: Introduce your little learners to the concepts and basics of the days of the week, months of the year, seasons, weather, number recognition, counting, skip counting, shapes, money, place value, and more with lots of hands-on activities and printable worksheets. Daily Math is a series of worksheets and drills created to reinforce 3rd-4th grade math skills. The skills covered include: - Basic Math Facts - Basic Skills: Place Value; Standard Measurements; Geometric Shape Identification; Ordinal Numbers; Number Patterns; Rounding Numbers; Writing Numbers as Words or Digits; Associative, Commutative, and Distributive Properties; Working with Zeroes; Order of Operations. - Intermediate Skills: Factors, Multiples, and Divisibility; Prime Numbers; Fractions and Percentages; Decimals; Subtracting Negative Numbers; Roman Numerals; Exponents; Value of Coins - Advanced Skills: Area, Perimeter, Circumference; Angles; Finding the Mean (Average); Variables - Miscellaneous: Adding Columns of Numbers; Solving for Rate; Comparing Ratios; Math Fun; Math Vocabulary; Probability; Describing/Defining Triangles; Inverse Operations - Word Problems (assorted) - Drills: Addition, Subtraction, Multiplication, and Division The Decimal Workshop is designed to help parents teach students through examples and practice problems. It is a great review unit for an elementary math class and can be done at the student’s own pace. Everyday Games: Over the years, Teresa has brought hundreds of games to our site that cover various reading and math skills. Here is a quick and handy way to access all the games by the specific skill they were created to reinforce. New games are added regularly. You can see all of the Everyday Games by clicking on the title above. Finally Conquer Fractions: Are you a non-math parent? Finally Conquer Fractions by Laura Baggett is designed to equip homeschool parents with the tools, confidence, and motivation needed for effectively teaching math to their children. This course focuses on the subject of fractions and provides a visual, step-by-step breakdown of what fractions are, how they work, and why they behave the way they do. Tips are also included for explaining fractions to your children. Fraction Workshop: This fun, twelve-week fraction workshop focuses on understanding the basics of fractions using a visual and hands-on approach. Designed for students of any age who need a boost in their understanding of fractions, this workshop will illustrate fractions using concrete objects—balls, vegetables, Popsicle sticks, and more. Geometry with Mr. D Math is a one-year video-based geometry class for high school. Each lesson includes video instruction with printable course work and detailed answer keys. Also included are twelve chapter tests and two semester exams. It is best suited for students who have an understanding of Algebra 1. Hands-On Math Help with Cuisenaire Rods is designed to empower parents to help their children learn. Its purpose is to help parents of students in preschool-3rd grade make math enjoyable instead of frustrating. Through a free app and weekly tutorial videos, Philip Rowlands shows parents how to teach math in a manner that involves play, directed activities, open-ended tasks, and challenges. It is designed to work with Cuisenaire rods, but instructions for making your own set are included if you do not currently have them. How to Teach Elementary Math: These videos from Dr. Price help demonstrate to parents various methods and strategies for teaching math to their elementary students. Worksheets for the students are also provided. New units are provided regularly. - Kindergarten: Rainbow Pairs to 10; Numbers 2 Before and After; Rainbow Facts Template; Easy Addition: Using a Near Ten Strategy to Add 8; Easy Subtraction: Near 10 Minus 9, Difference of 9; Addition: Doubles Facts - Grade One: Count on 2 within 10 Addition Facts; Double +1 Addition Facts; Addition & Subtraction Revision: Difference of 1, 2, or 3; Count On 1 Addition Facts; Count On 1, Count Back 1; Ten Frame Introduction; Near 10; Adding 8; Count on 2 Addition Facts; Count on 1, 2, or 3 Addition Facts; Ten Frame Flash Cards; Subtracting 0 and 10; Counting On and Back by 1; Addition Fact Families: Counting On and Back; Subtraction from 2-Digit Numbers Using a Written Algorithm; Checkup Test: Easy Multiplication 5x, 3x; Probability: Outcomes of Familiar Events; Easy Multiplication: 3x Tables - Grade Two: Count Back 3 Subtraction Facts; Multiplying by 10; Addition Number Fact Families; Converting Fractions (Easy Level); Counting Fractions Along a Number Line; Introducing Fractions; Common Fractions: Halves, Fourths, and Quarters; Addition & Subtraction Revision; Doubling and Halving; Multiplication and Division Fact Families; Revision of Easy Multiplication and Division Facts; Counting Fractions Along a Number Line; Easy Division; Doubles Addition Facts; Subtraction: Halving; Adding 9 - Grade Three: Multiplication and Division Revision Facts x 3; Near 10 Facts: Adding and Subtracting 8; Addition & Subtraction Revision; Multiplication and Division Revision; Equivalent Fractions; Addition & Subtraction Rainbow Facts; Double Double Strategy: 4x Facts; 9x Tables; Dividing by 5 and by 10; Dividing by 4; Roman Numerals to 399; 9x Tables and Strategies; Multiplication Thinking Strategies Poster; Multiplication Tables Revision; Hundreds, Tens & Ones Place Value: Base Ten Blocks; All Operations Revision: Remaining Facts, 6x - Grade Four: Division with Remainders; Revision of Multiplication and Division by 8; Divisibility by 3 and 9; Extended Multiplication and Division: 3x; Rainbow Facts to 100; 7x Multiplication Facts Revision; Division by 5 & 10 with Remainders; Division with Remainders: Dividing by 2 and 4; Decimal Fractions: Skip Counting in Hundredths; Comparing Fractions; Mental Strategies: Adding Nice Numbers; Factors & Multiples: Lowest Common Multiple; Mental Strategies: Subtracting 100 - Grade Five: Adding “Nice” Numbers; Dividing by 10, 100, and 1,000; Adding and Subtracting Fractions with Common Denominators; Adding Fractions with Like Denominators (Harder Level); Simplifying Fractions; Converting Between Improper Fractions and Mixed Numbers; Converting Common Fractions into Decimals; Doubling Large Numbers; Calculating 1% and 0.5% of an Amount; Revision: Order of Operations; Difference of; Multiplying and Dividing by 5 and 10; Factors and Multiples of 2 and 4; Probability with Spinners; Converting Common Fractions to Decimals - Why Teach Math and Developing Number Fluency - Why Teach Math - Mental Strategies - Development of Mental Strategies - The Dream of Students’ Success in Math - Using Physical and Visual Models for Numbers - Helping Children to Develop Math Thinking - Developing “Automatic” Responses in Number Facts - A Four-Phase Model for Teaching Operations - More Ten Frames Examples - Modeling Multiplication with Arrays - Developing Confidence and Freedom via Thinking Strategies - A Systematic Approach, Not a Short-Term Fix - Show ANY Number with a Number Line - Combining Math and Old Technology - Using Public Spaces to Teach Symmetry - How to teach Multiplication Facts I - How to teach Multiplication Facts II - Math and Critical Thinking - What is REALLY Important in Math? - Teaching Slope in the Mountains of Switzerland - Five Misconceptions About the Teaching of Math to Young Children - Five Benefits of Memorizing Your Times Tables - Six Simple, Free Activities any Parent can use at Home to Help Their Child With Math - Sundials and Math in Switzerland The Multiplication and Division Practice Unit is designed to help parents teach students multiplication and division through examples and practice problems. It is a great review unit for an elementary math class and can be done at the student’s own pace. The Multiplication Workshop is a step-by-step guide to learning or practicing multiplication tables using written problems, manipulatives, flash cards, word problems, and speed drills. The 160 worksheets spend the most time working on the concept of multiplication and learning the 1-12 times tables. As a bonus at the end of the workshop, multiplication with multiple-digit numbers is introduced. Pre-Algebra: This is a full forty-week pre-algebra course that includes instruction on fractions, decimals, percents, English and Metric measures, number sets and graphing, algebraic expressions, solving algebraic equations, geometry, statistics ,and probability. Precalculus: This 36-week course is designed for students who have completed Algebra 1, Algebra 2, and trigonometry. This course covers relations and functions, linear and quadratic functions, polynomial and rational functions, exponential and logarithmic functions, conics, systems of equations, matrices, sequences, and binomial theorem. It concludes with a final test. Special Needs: Special Needs Educational Consultant, Judi Munday, has put together a 19-part program to help those with children with special needs, be it language delays or deficits, dyslexia, or difficulty with math. She provides practical solutions and helpful tips for choosing curriculum and modifying and adapting instruction. An explanation of the skills of learning—acquiring, organizing, storing, and retrieving information—and where possible breakdowns may occur along the path is helpful for those who are seeking to understand the best course of action to facilitate their child’s success. The following list highlights the lessons specific to language-based or reading challenges. You can see the complete list of Judi’s articles here. - Lesson Fourteen: Assistive Technology: Struggling with Math Starting Out with First Grade Math is a full-year course that introduces first graders to the basic math concepts they need to build a strong foundation. Through interactive slides and printable assignments, students explore numbers, greater and less than, addition, subtraction, fact families, and much more. Statistics is a six-month course that covers how statistics are used, vocabulary, normal distribution, graphs, scatter diagrams, pie charts, standard deviation, the “68-95-99.7 rule,” slope, uses of slope, equation of a line, demographics, sample, sample size, the use of statistics in presidential politics and in TV ratings, and more. Stepping Up with Second Grade Math is a full-year course that builds on math knowledge previously learned to help second graders learn additional key math skills including skip counting, place value, regrouping, story problems, multiplication, division, geometric vocabulary, and much more. Stretching Higher with Third Grade Math builds on the math skills students have previously learned and guides them through new concepts involving multiplication, division, working with money, fractions, decimals, equations, and more. Daily lessons are provided with ample problems for review and practice. Regular quizzes and tests are also included. Answer keys are provided. Algebra for Kids Supplemental Worksheets-Skip Counting and Multiples: Worksheets by Bob Hazen provide thirteen supplemental math worksheets that enable young students to practice skills such as skip counting and multiples. The Trigonometry course on SchoolhouseTeachers.com is presented as a supplemental trigonometry course. The topics cover the ideas and topics that are basic in working with trigonometry including right angles, acute angles, graphing, inverse functions, solving trigonometric equations, and more. The Understanding Fractions Workshop is designed to help parents teach students fraction skills through examples and practice problems. It is a great review unit for an elementary math class and can be done at the student’s own pace. The Whole Number Place Value Workshop is designed to help parents teach students place value through examples and practice problems. It is a great review unit for an elementary math class and can be done at the student’s own pace.
The Discrete Fourier Transform converts discrete data from a time wave into a frequency spectrum. Using the DFT implies that the finite segment that is analyzed is one period of an infinitely extended periodic signal. The DFT equation: x(k) is the time wave that is converted to a frequency spectrum by the DFT. Here are key concepts required to understand a DFT: - The "sampling rate", sr. The sampling rate is the number of samples taken over a time period. For simplicity we will make the time interval between samples equal. This is the "sample interval", si. - The fundamental period, T, is the period of all the samples taken. This is also called the "window". - The "fundamental frequency" is f0, which is 1/T. f0 is the first harmonic, the second harmonic is 2f0, the third is 3f0, etc. - The number of samples is N. - The "Nyquist Frequency", fc, is half the sampling rate. The Nyquist frequency is the maximum frequency that can be detected for a given sampling rate. This is because in order to measure a wave you need at least two sample points to identify it (trough and peak). - "Euler's formula" -- - The sampled part of the time wave, x(t), should be "typical" of how the wave behaves over all time that it exists. - This notation makes handling the exponential easier. This is sometimes called the "twiddle factor." For simplicity, we will sample a sine wave with a small number of points, N, and perform a DFT on it, then we will employ each of the concepts above. Note, the sine wave is a time wave, and could be any wave in nature, for example a sound wave. The horizontal axis is time. The vertical axis is amplitude. Notice how in the diagram above we are sampling four points. The fundamental period, T, of the wave sampled is set to 2pi. This applies to any wave we want to sample. The interval between samples is 2pi/N, so in this case it is 2pi/4. Thus, the interval between samples is pi/2 in this case. The time wave is thus, x(k) = sin(pi/2k) for k = 0 to N -1. The last point sampled is always the point just before 2pi, because the wave is considered to be a repeating pattern and wraps around back to the value at k = 0, so you aren't missing any information. We also need to know the time taken to sample the wave, so that we can tie it to a frequency. In our example, the time taken for the fundamental period, T, is 0.1 seconds (this value is measured when the wave is captured). That means the sine wave is a 10 Hz wave. Hertz = cycles per second. Also, the sampling interval, si, is the fundamental period time divided by the number of samples. So, si = T/N = (0.1)/4 seconds, or 0.025 seconds. The sampling rate, or frequency, sr, = 1/si = 40 Hz, or 40 samples per second. For the sine wave, the value at each of the four points sampled is: And, before we plug into the DFT, some more on Wn, the twiddle factor, referenced above: The DFT formula, then, for a four point sample and with the twiddle factor is: Now, Euler's Formula for N=4: For the equation above, where kn = 0 to N - 1, i.e. 0 to 3, here are the results: Notice that any additional integer values of kn will cycle back around. For example, kn = 4 cycles back to kn=0, so the value is 1. kn = 5 cycles back around to kn = 1, so the value is -j. The equation "kn modulus 4" determines which value of W is selected. Also, note that for larger samples the cycle is bigger. So for N=8 the equation would be "kn modulus 8". This is probably why W is called the "twiddle factor". Now put this together for the DFT: Here is the DFT worked out for all four points and for four frequencies: Evaluating the output data. Each F(n) value outputs a phase at a particular frequency. The frequency of the point is determined by the fundamental frequency multiplied by n, i.e. f = f0n, where f0=1/T = 10Hz. The output values are the phase of the frequencies, which are represented by a real part and an imaginary part thus: real + jimaginary. The fundamental frequency, first harmonic, is 10 Hz as calculated above. The magnitude at a frequency is Calculated thus sqrt(realreal + imaginaryimaginary). Below is a frequency spectrum plot for the sine wave determined from the DFT we just worked through: The frequency plot is in the "frequency domain". The magnitudes are plotted in Diagram 2. The spike at 10 Hz shows that the DFT pulled out one of the frequencies that is in the sine wave. In fact, the sine wave is a 10 Hz sine wave, so that makes sense. However, the spike at 30 Hz should not be there, because there is no 30 Hz wave in the sine wave. So what accounts for that spike? Well, this is where the Nyquist Frequency, fc, mentioned above comes in. The Nyquist frequency is the cut off point above which the data from the DFT is no longer valid. The sampling rate is 40 Hz, and fc is half the sampling frequency, which means that any frequency above 20 Hz will not be valid in this case. So, the 30 Hz frequency is a spurious signal. That completes analysis of a very simple wave. Most waves will have many more frequencies in them, and thus many more spikes of various magnitudes along the frequency spectrum. For example, Diagram 3, below, is a plot of a triangle wave in time and its corresponding frequency spectrum: The Next section is on the FFT. The FFT builds on the knowledge of the DFT described above, so it should be understood before moving on to the FFT. \|<< Previous\|\|Next >>\|
You have been given nine weights, one of which is slightly heavier than the rest. Can you work out which weight is heavier in just two weighings of the balance? The NRICH team are always looking for new ways to engage teachers and pupils in problem solving. Here we explain the thinking behind Can you arrange the numbers 1 to 17 in a row so that each adjacent pair adds up to a square number? Do you notice anything about the solutions when you add and/or subtract consecutive negative numbers? The number of plants in Mr McGregor's magic potting shed increases overnight. He'd like to put the same number of plants in each of his gardens, planting one garden each day. How can he do it? Advent Calendar 2011 - a mathematical activity for each day during the run-up to Christmas. Can you find six numbers to go in the Daisy from which you can make all the numbers from 1 to a number bigger than 25? An investigation involving adding and subtracting sets of consecutive numbers. Lots to find out, lots to explore. Bellringers have a special way to write down the patterns they ring. Learn about these patterns and draw some of your own. Try to solve this very difficult problem and then study our two suggested solutions. How would you use your knowledge to try to solve variants on the original problem? Different combinations of the weights available allow you to make different totals. Which totals can you make? Charlie and Abi put a counter on 42. They wondered if they could visit all the other numbers on their 1-100 board, moving the counter using just these two operations: x2 and -5. What do you think? Many numbers can be expressed as the sum of two or more consecutive integers. For example, 15=7+8 and 10=1+2+3+4. Can you say which numbers can be expressed in this way? You have twelve weights, one of which is different from the rest. Using just 3 weighings, can you identify which weight is the odd one out, and whether it is heavier or lighter than the rest? Use the interactivity to listen to the bells ringing a pattern. Now it's your turn! Play one of the bells yourself. How do you know when it is your turn to ring? A 2 by 3 rectangle contains 8 squares and a 3 by 4 rectangle contains 20 squares. What size rectangle(s) contain(s) exactly 100 squares? Can you find them all? How many solutions can you find to this sum? Each of the different letters stands for a different number. This article for teachers describes several games, found on the site, all of which have a related structure that can be used to develop the skills of strategic planning. Find out about Magic Squares in this article written for students. Why are they magic?! My two digit number is special because adding the sum of its digits to the product of its digits gives me my original number. What could my number be? Problem solving is at the heart of the NRICH site. All the problems give learners opportunities to learn, develop or use mathematical concepts and skills. Read here for more information. Ben passed a third of his counters to Jack, Jack passed a quarter of his counters to Emma and Emma passed a fifth of her counters to Ben. After this they all had the same number of counters. An irregular tetrahedron is composed of four different triangles. Can such a tetrahedron be constructed where the side lengths are 4, 5, 6, 7, 8 and 9 units of length? A game for 2 people. Take turns placing a counter on the star. You win when you have completed a line of 3 in your colour. Label this plum tree graph to make it totally magic! Use the interactivity to play two of the bells in a pattern. How do you know when it is your turn to ring, and how do you know which bell to ring? First Connect Three game for an adult and child. Use the dice numbers and either addition or subtraction to get three numbers in a straight line. A Latin square of order n is an array of n symbols in which each symbol occurs exactly once in each row and exactly once in each column. Use the differences to find the solution to this Sudoku. A Sudoku with a twist. This challenging activity involves finding different ways to distribute fifteen items among four sets, when the sets must include three, four, five and six items. This tricky challenge asks you to find ways of going across rectangles, going through exactly ten squares. A challenging activity focusing on finding all possible ways of stacking rods. Given the products of adjacent cells, can you complete this Sudoku? Place the 16 different combinations of cup/saucer in this 4 by 4 arrangement so that no row or column contains more than one cup or saucer of the same colour. The clues for this Sudoku are the product of the numbers in adjacent squares. There are nine teddies in Teddy Town - three red, three blue and three yellow. There are also nine houses, three of each colour. Can you put them on the map of Teddy Town according to the rules? Four friends must cross a bridge. How can they all cross it in just 17 minutes? If you have only 40 metres of fencing available, what is the maximum area of land you can fence off? Rather than using the numbers 1-9, this sudoku uses the nine different letters used to make the words "Advent Calendar". A package contains a set of resources designed to develop students’ mathematical thinking. This package places a particular emphasis on “being systematic” and is designed to meet. . . . You are given the Lowest Common Multiples of sets of digits. Find the digits and then solve the Sudoku. This cube has ink on each face which leaves marks on paper as it is rolled. Can you work out what is on each face and the route it has taken? Make your own double-sided magic square. But can you complete both sides once you've made the pieces? If you take a three by three square on a 1-10 addition square and multiply the diagonally opposite numbers together, what is the difference between these products. Why? Take three whole numbers. The differences between them give you three new numbers. Find the differences between the new numbers and keep repeating this. What happens? This is a variation of sudoku which contains a set of special clue-numbers. Each set of 4 small digits stands for the numbers in the four cells of the grid adjacent to this set. An extra constraint means this Sudoku requires you to think in diagonals as well as horizontal and vertical lines and boxes of This challenge extends the Plants investigation so now four or more children are involved. How many different symmetrical shapes can you make by shading triangles or squares?
Solar time is a reckoning of the passage of time based on the Sun's position in the sky. The fundamental unit of solar time is the day. Two types of solar time are apparent solar time (sundial time) and mean solar time (clock time). Fix a tall pole vertically in the ground; at some instant on any sunny day the shadow will point exactly north or south (or disappear, if the Sun is directly overhead). That instant is local apparent noon: 12:00 local apparent time. About 24 hours later the shadow will again point north/south, the Sun seeming to have covered a 360-degree arc around the Earth's axis. When the Sun has covered exactly 15 degrees (1/24 of a circle, both angles being measured in a plane perpendicular to the Earth's axis), local apparent time is 13:00 exactly; after 15 more degrees it will be 14:00 exactly. The problem is that in September the Sun takes less time (as measured by an accurate clock) to make an apparent revolution than it does in December; 24 "hours" of solar time can be 21 seconds less or 29 seconds more than 24 hours of clock time. As explained in the equation of time article, this is due to the eccentricity of the Earth's orbit (i.e. the Earth's orbit is not perfectly circular, meaning that the Earth-Sun distance varies throughout the year), and the fact that the Earth's axis is not perpendicular to the plane of its orbit (the so-called obliquity of the ecliptic). The effect of this is that a clock running at a constant rate – e.g. completing the same number of pendulum swings in each hour – cannot follow the actual Sun; instead it follows an imaginary "mean Sun" that moves along the celestial equator at a constant rate that matches the real Sun's average rate over the year. This is "mean solar time", which is still not perfectly constant from one century to the next but is close enough for most purposes. Currently a mean solar day is about 86,400.002 SI seconds. The two kinds of solar time (apparent solar time and mean solar time) are among the three kinds of time reckoning that were employed by astronomers until the 1950s. (The third kind of traditional time reckoning is sidereal time, which is based on the apparent motions of stars other than the Sun.) By the 1950s it had become clear that the Earth's rotation rate was not constant, so astronomers developed ephemeris time, a time scale based on the positions of solar system bodies in their orbits. The apparent sun is the true sun as seen by an observer on Earth. Apparent solar time or true solar time is based on the apparent motion of the actual Sun. It is based on the apparent solar day, the interval between two successive returns of the Sun to the local meridian. Solar time can be crudely measured by a sundial. The equivalent on other planets is termed local true solar time (LTST). The length of a solar day varies through the year, and the accumulated effect produces seasonal deviations of up to 16 minutes from the mean. The effect has two main causes. First, Earth's orbit is an ellipse, not a circle, so the Earth moves faster when it is nearest the Sun (perihelion) and slower when it is farthest from the Sun (aphelion) (see Kepler's laws of planetary motion). Second, due to Earth's axial tilt (known as the obliquity of the ecliptic), the Sun's annual motion is along a great circle (the ecliptic) that is tilted to Earth's celestial equator. When the Sun crosses the equator at both equinoxes, the Sun's daily shift (relative to the background stars) is at an angle to the equator, so the projection of this shift onto the equator is less than its average for the year; when the Sun is farthest from the equator at both solstices, the Sun's shift in position from one day to the next is parallel to the equator, so the projection onto the equator of this shift is larger than the average for the year (see tropical year). In June and December when the sun is farthest from the celestial equator a given shift along the ecliptic corresponds to a large shift at the equator. So apparent solar days are shorter in March and September than in June or December. \|Date\|\|Duration in mean solar time\| \|February 11\|\|24 hours\| \|March 26\|\|24 hours − 18.1 seconds\| \|May 14\|\|24 hours\| \|June 19\|\|24 hours + 13.1 seconds\| \|July 25/26\|\|24 hours\| \|September 16\|\|24 hours − 21.3 seconds\| \|November 2/3\|\|24 hours\| \|December 22\|\|24 hours + 29.9 seconds\| These lengths will change slightly in a few years and significantly in thousands of years. Mean solar time is the hour angle of the mean Sun plus 12 hours. Currently (2009) this is realized with the UT1 time scale, constructed mathematically from very long baseline interferometry observations of the diurnal motions of radio sources located in other galaxies, and other observations. The duration of daylight varies during the year but the length of a mean solar day is nearly constant, unlike that of an apparent solar day. An apparent solar day can be 20 seconds shorter or 30 seconds longer than a mean solar day. Long or short days occur in succession, so the difference builds up until mean time is ahead of apparent time by about 14 minutes near February 6 and behind apparent time by about 16 minutes near November 3. The equation of time is this difference, which is cyclical and does not accumulate from year to year. Mean time follows the "mean sun", best described by Meeus: The length of the mean solar day is slowly increasing due to the tidal acceleration of the Moon by the Earth and the corresponding slowing of Earth's rotation by the Moon. Many methods have been used to simulate mean solar time. The earliest were clepsydras or water clocks, used for almost four millennia from as early as the middle of the 2nd millennium BC until the early 2nd millennium. Before the middle of the 1st millennium BC, the water clocks were only adjusted to agree with the apparent solar day, thus were no better than the shadow cast by a gnomon (a vertical pole), except that they could be used at night. But it has long been known that the Sun moves eastward relative to the fixed stars along the ecliptic. Since the middle of the first millennium BC the diurnal rotation of the fixed stars has been used to determine mean solar time, against which clocks were compared to determine their error rate. Babylonian astronomers knew of the equation of time and were correcting for it as well as the different rotation rate of the stars, sidereal time, to obtain a mean solar time much more accurate than their water clocks. This ideal mean solar time has been used ever since then to describe the motions of the planets, Moon, and Sun. Mechanical clocks did not achieve the accuracy of Earth's "star clock" until the beginning of the 20th century. Today's atomic clocks have a much more constant rate than the Earth, but its star clock is still used to determine mean solar time. Since sometime in the late 20th century, Earth's rotation has been defined relative to an ensemble of extra-galactic radio sources and then converted to mean solar time by an adopted ratio. The difference between this calculated mean solar time and Coordinated Universal Time (UTC) determines whether a leap second is needed. (The UTC time scale now runs on SI seconds, and the SI second, when adopted, was already a little shorter than the current value of the second of mean solar time.)
All living beings have genes. They exist throughout the body. Genes are a set of instructions that determine what the organism is like, its appearance, how it survives, and how it behaves in its environment. Genes are made of a substance called deoxyribonucleic acid, or DNA. They give instructions for a living being to make molecules called proteins. A geneticist is a person who studies genes and how they can be targeted to improve aspects of life. Genetic engineering can provide a range of benefits for people, for example, increasing the productivity of food plants or preventing diseases in humans. Genes are a section of DNA that are in charge of different functions like making proteins. Long strands of DNA with lots of genes make up chromosomes. DNA molecules are found in chromosomes. Chromosomes are located inside of the nucleus of cells. Each chromosome is one long single molecule of DNA. This DNA contains important genetic information. Genes vary in complexity. In humans, they range in size from a few hundred DNA bases to more than 2 million bases. Different living things have different shapes and numbers of chromosomes. Humans have DNA contains the biological instructions that make each species unique. DNA is passed from adult organisms to their offspring during reproduction. The building blocks of DNA are called nucleotides. Nucleotides have three parts: A phosphate group, a sugar group and one of four types of nitrogen bases. A gene consists of a long combination of four different nucleotide bases, or chemicals. There are many possible combinations. The four nucleotides are: - A (adenine) - C (cytosine) - G (guanine) - T (thymine) To recap in more detail: Genes carry the codes ACGT. Each person has thousands of genes. They are like a computer program, and they make the individual what they are. A gene is a tiny section of a long DNA double helix molecule, which consists of a linear sequence of base pairs. A gene is any section along the DNA with instructions encoded that allow a cell to produce a specific product – usually a protein, such as an enzyme – that triggers one precise action. DNA is the chemical that appears in strands. Every cell in a person’s body has the same DNA, but each person’s DNA is different. This is what makes each person unique. DNA is made up of two long-paired strands spiraled into the famous double helix. Each strand Genes decide almost everything about a living being. One or more genes can affect a specific trait. Genes may interact with an individual’s environment too and change what the gene makes. Genes affect hundreds of internal and external factors, such as whether a person will get a particular color of eyes or what diseases they may develop. A gene is a basic unit of heredity in a living organism. Genes come from our parents. We may inherit our physical traits and the likelihood of getting certain diseases and conditions from a parent. Genes contain the data needed to build and maintain cells and pass genetic information to offspring. Each cell contains two sets of chromosomes: One set comes from the mother and the other comes from the father. The male sperm and the female egg carry a single set of 23 chromosomes each, including 22 autosomes plus an X or Y sex chromosome. A female inherits an X chromosome from each parent, but a male inherits an X chromosome from their mother and a Y chromosome from their father. It aims to determine the sequence of the chemical pairs that make up human DNA and to identify and map the 20,000 to 25,000 or so genes that make up the human genome. The project was started in 1990 by a group of international researchers, the United States’ National Institutes of Health (NIH) and the Department of Energy. The goal was to sequence 3 billion letters, or base pairs, in the human genome, that make up the complete set of DNA in the human body. By doing this, the scientists hoped to provide researchers with powerful tools, not only to understand the genetic factors in human disease, but also to open the door for new strategies for diagnosis, treatment, and prevention. The HGP was completed in 2003, and all the data generated is available for free access on the internet. Apart from humans, the HGP also looked at other organisms and animals, such as the fruit fly and E. coli. Over three billion nucleotide combinations, or combinations of ACGT, have been found in the human genome, or the collection of genetic features that can make up the human body. Mapping the human genome brings scientists closer to developing effective treatments for hundreds of diseases. The project has fueled the discovery of more than 1,800 disease genes. This has made it easier for researchers to find a gene that is suspected of causing an inherited disease in a matter of days. Before this research was carried out, it could have taken years to find the gene. Genetic tests can show an individual whether they have a genetic risk for a specific disease. The results can help healthcare professionals diagnose conditions. The HGP is expected to speed up progress in medicine, but there is still much to learn, especially regarding how genes behave and how they can be used in treatment. At least 350 biotechnology-based products are currently in clinical trials. In 2005, the HapMap, a catalog of common genetic variation or haplotypes in the human genome, was created. This data has helped to speed up the search for the genes involved in common human diseases. In recent years, geneticists have found another layer of heritable genetic data that is not held in the genome, but in the “epigenome,” a group of chemical compounds that can tell the genome what to do. In the body, DNA holds the instructions for building proteins, and these proteins are responsible for a number of functions in a cell. The epigenome is made up of chemical compounds and proteins that can attach to DNA and direct a variety of actions. These actions include turning genes on and off. This can control the production of proteins in particular cells. Gene switches can turn genes on and off at different times and for different lengths of time. Recently, scientists have discovered genetic switches that increase the lifespan and boost fitness in worms. They believe these could be linked to an increased lifespan in mammals. The genetic switches that they have discovered involve enzymes that are ramped up after mild stress during early development. This increase in enzyme production continues to affect the expression of genes throughout the animal’s life. This could lead to a breakthrough in the goal to develop drugs that can flip these switches to improve human metabolic function and increase longevity. When epigenomic compounds attach themselves to DNA in the cell and modify the function, they are said to have “marked” the genome. The marks do not change the sequence of the DNA, but they do change the way cells use the DNA’s instructions. The marks can be passed on from cell to cell as they divide, and they can even be passed from one generation to the next. Specialized cells can The chemical tags on the DNA and histones can become rearranged as the specialized cells and the epigenome change throughout a person’s lifetime. Lifestyle and environmental factors such as smoking, diet and infectious diseases can bring about changes in the epigenome. They can expose a person to pressures that prompt chemical responses. These responses can lead to direct changes in the epigenome, and some of these changes Some of these changes are linked to the development of disease. Cancer can result from changes in the genome, the epigenome or both. Changes in the epigenome can switch on or off the genes that are involved in cell growth or the immune response. These changes can cause uncontrolled growth, a feature of cancer, or a failure of the immune system to destroy tumors. Researchers in The Cancer Genome Atlas (TCGA) network are comparing the genomes and epigenomes of normal cells with those of cancer cells in the hope of compiling a current and complete list of possible epigenomic changes that can lead to cancer. Researchers in epigenomics are focused on trying to In gene therapy, genes are inserted into a patient’s cells and tissues to treat a disease, usually a hereditary disease. Gene therapy uses sections of DNA to treat or prevent disease. This science is still in its early stages, but there has been some success. For example, in 2016, scientists reported that they had managed to improve the eyesight of 3 adult patients with congenital blindness by using gene therapy. In 2017, a reproductive endocrinologist, named John Zhang, and a team at the New Hope Fertility Center in New York used a technique called mitochondrial replacement therapy in a revolutionary way. They announced the birth of a child to a mother carrying a fatal genetic defect. Researchers combined DNA from two women and one man to bypass the defect. The result was a healthy baby boy with three genetic parents. This type of research is still in the early stages, and much is still unknown, but results look promising. Scientists are looking at different ways of treating cancer using gene therapy. Experimental gene therapy may use patients’ own blood cells to kills cancer cells. In one study, 82 percent of patients had their cancer shrink by at least half at some point during treatment. Gene testing to predict cancer Women with the BRCA1 gene have a significantly higher chance of developing breast cancer. A woman can have a test to find out whether she carries that gene. BRCA1 carriers have a Genetic tests for personalized therapy Scientists say that one day we will be able to test a patient to find out which specific medicines are best for them, depending on their genetic makeup. Some medicines work well for some patients, but not for others. Gene therapy is still a growing science, but in time, it may become a viable medical treatment.
Learning Goal: To learn the quantitative use of the lens equation, as well as how to determine qualitative properties of solutions. In working with lenses, there are three important quantities to consider: The object distance s is the distance along the axis of the lens to the object. The image distance s' is the distance along the axis of the lens to the image. The focal length f is an intrinsic property of the lens. These three quantities are related through the equation 1/s + 1/s' = 1/f. Note that this equation is valid only for thin, spherical lenses. Unless otherwise specified, a lens problem always assumes that you are using thin, spherical lenses. The equation above allows you to calculate the locations of images and objects. Frequently, you will also be interested in the size of the image or object, particularly if you are considering a magnifying glass or microscope. The ratio of the size of an image to the size of the object is called the magnification. It is given by: m = y'/y = -s'/s, where y' is the height of the image and y is the height of the object. The second equality allows you to find the size of the image (or object) with the information provided by the thin lens equation. All of the quantities in the above equations can take both positive and negative values. Positive distances correspond to real images or objects, while negative distances correspond to virtual images or objects. Positive heights correspond to upright images or objects, while negative heights correspond to inverted images or objects. The following table summarizes these properties: s real virtual s' real virtual y upright inverted y' upright inverted The focal length f can also be positive or negative. A positive focal length corresponds to a converging lens, while a negative focal length corresponds to a diverging lens. Consider an object with s = 12 cm that produces an image with s' = 15 cm. Note that whenever you are working with a physical object, the object distance will be positive (in multiple optics setups, you will encounter "objects" that are actually images, but that is not a possibility in this problem). A positive image distance means that the image is formed on the side of the lens from which the light emerges. Now consider a diverging lens with focal length f=-15 cm, producing an upright image that is 5/9 as tall as the object. f) What is the object distance? g) What is the image distance? A lens placed at the origin with its axis pointing along the x axis produces a real inverted image at x=-24 cm that is twice as tall as the object. h) What is the image distance? i) What is the x coordinate of the object? Keep in mind that a real image and a real object should be on opposite sides of the lens. j) The lens is converging k) What is the focal length of the lens in cm? Please see attachment for full solutions. f) You will need to use the magnification equation to find a relationship between s and s'. Then substitute into the thin lens equation to solve for s. The answer is s = 12 cm. g) A lens placed at the origin with its ... It shows how to apply thin lens equation to find the object, image, and magnification. Solution includes .pdf and word attachments.
The octal numeral system, or oct for short, is the base-8 number system, and uses the digits 0 to 7. Octal numerals can be made from binary numerals by grouping consecutive binary digits into groups of three (starting from the right). For example, the binary representation for decimal 74 is 1001010. Two zeroes can be added at the left: (00)1 001 010, corresponding the octal digits 1 1 2, yielding the octal representation 112. In the decimal system each decimal place is a power of ten. For example: In the octal system each place is a power of eight. For example: By performing the calculation above in the familiar decimal system we see why 112 in octal is equal to 64+8+2 = 74 in decimal. By Native Americans - It has been suggested that the reconstructed Proto-Indo-European word for "nine" might be related to the PIE word for "new". Based on this, some have speculated that proto-Indo-Europeans used an octal number system, though the evidence supporting this is slim. - In 1668, John Wilkins in An Essay towards a Real Character, and a Philosophical Language proposed use of base 8 instead of 10 "because the way of Dichotomy or Bipartition being the most natural and easie kind of Division, that Number is capable of this down to an Unite". - In 1716, King Charles XII of Sweden asked Emanuel Swedenborg to elaborate a number system based on 64 instead of 10. Swedenborg however argued that for people with less intelligence than the king such a big base would be too difficult and instead proposed 8 as the base. In 1718 Swedenborg wrote (but did not publish) a manuscript: "En ny rekenkonst som om vexlas wid Thalet 8 i stelle then wanliga wid Thalet 10" ("A new arithmetic (or art of counting) which changes at the Number 8 instead of the usual at the Number 10"). The numbers 1-7 are there denoted by the consonants l, s, n, m, t, f, u (v) and zero by the vowel o. Thus 8 = "lo", 16 = "so", 24 = "no", 64 = "loo", 512 = "looo" etc. Numbers with consecutive consonants are pronounced with vowel sounds between in accordance with a special rule. - Writing under the pseudonym "Hirossa Ap-Iccim" in The Gentleman's Magazine, (London) July 1745, Hugh Jones proposed an octal system for British coins, weights and measures. "Whereas reason and convenience indicate to us an uniform standard for all quantities; which I shall call the Georgian standard; and that is only to divide every integer in each species into eight equal parts, and every part again into 8 real or imaginary particles, as far as is necessary. For tho' all nations count universally by tens (originally occasioned by the number of digits on both hands) yet 8 is a far more complete and commodious number; since it is divisible into halves, quarters, and half quarters (or units) without a fraction, of which subdivision ten is uncapable...." In a later treatise on Octave computation (1753) Jones concluded: "Arithmetic by Octaves seems most agreeable to the Nature of Things, and therefore may be called Natural Arithmetic in Opposition to that now in Use, by Decades; which may be esteemed Artificial Arithmetic." - In 1801, James Anderson criticized the French for basing the metric system on decimal arithmetic. He suggested base 8, for which he coined the term octal. His work was intended as recreational mathematics, but he suggested a purely octal system of weights and measures and observed that the existing system of English units was already, to a remarkable extent, an octal system. - In the mid-19th century, Alfred B. Taylor concluded that "Our octonary [base 8] radix is, therefore, beyond all comparison the "best possible one" for an arithmetical system." The proposal included a graphical notation for the digits and new names for the numbers, suggesting that we should count "un, du, the, fo, pa, se, ki, unty, unty-un, unty-du" and so on, with successive multiples of eight named "unty, duty, thety, foty, paty, sety, kity and under." So, for example, the number 65 (101 in octal) would be spoken in octonary as under-un. Taylor also republished some of Swedenborg's work on octal as an appendix to the above-cited publications. Octal became widely used in computing when systems such as the UNIVAC 1050, PDP-8, ICL 1900 and IBM mainframes employed 6-bit, 12-bit, 24-bit or 36-bit words. Octal was an ideal abbreviation of binary for these machines because their word size is divisible by three (each octal digit represents three binary digits). So two, four, eight or twelve digits could concisely display an entire machine word. It also cut costs by allowing Nixie tubes, seven-segment displays, and calculators to be used for the operator consoles, where binary displays were too complex to use, decimal displays needed complex hardware to convert radices, and hexadecimal displays needed to display more numerals. All modern computing platforms, however, use 16-, 32-, or 64-bit words, further divided into eight-bit bytes. On such systems three octal digits per byte would be required, with the most significant octal digit representing two binary digits (plus one bit of the next significant byte, if any). Octal representation of a 16-bit word requires 6 digits, but the most significant octal digit represents (quite inelegantly) only one bit (0 or 1). This representation offers no way to easily read the most significant byte, because it's smeared over four octal digits. Therefore, hexadecimal is more commonly used in programming languages today, since two hexadecimal digits exactly specify one byte. Some platforms with a power-of-two word size still have instruction subwords that are more easily understood if displayed in octal; this includes the PDP-11 and Motorola 68000 family. The modern-day ubiquitous x86 architecture belongs to this category as well, but octal is rarely used on this platform, although certain properties of the binary encoding of opcodes become more readily apparent when displayed in octal, e.g. the ModRM byte, which is divided into fields of 2, 3, and 3 bits, so octal can be useful in describing these encodings. Before the availability of assemblers, some programmers would handcode programs in octal; for instance, Dick Whipple and John Arnold wrote Tiny BASIC Extended directly in machine code, using octal. Octal is sometimes used in computing instead of hexadecimal, perhaps most often in modern times in conjunction with file permissions under Unix systems (see chmod). It has the advantage of not requiring any extra symbols as digits (the hexadecimal system is base-16 and therefore needs six additional symbols beyond 0–9). It is also used for digital displays. In programming languages, octal literals are typically identified with a variety of prefixes, including the digit 0, the letters o or q, the digit–letter combination 0o, or the symbol & or $. In Motorola convention, octal numbers are prefixed with @, whereas a small (or capital) letter o or q is added as a postfix following the Intel convention. In Concurrent DOS, Multiuser DOS and REAL/32 as well as in DOS Plus and DR-DOS various environment variables like $CLS, $ON, $OFF, $HEADER or $FOOTER support an \nnn octal number notation, and DR-DOS DEBUG utilizes \ to prefix octal numbers as well. For example, the literal 73 (base 8) might be represented as 073, o73, q73, 0o73, \73, @73, &73, $73 or 73o in various languages. Octal numbers that are used in some programming languages (C, Perl, PostScript…) for textual/graphical representations of byte strings when some byte values (unrepresented in a code page, non-graphical, having special meaning in current context or otherwise undesired) have to be to escaped as \nnn. Octal representation may be particularly handy with non-ASCII bytes of UTF-8, which encodes groups of 6 bits, and where any start byte has octal value \3nn and any continuation byte has octal value \2nn. Conversion between bases Decimal to octal conversion Method of successive Euclidean division by 8 To convert integer decimals to octal, divide the original number by the largest possible power of 8 and divide the remainders by successively smaller powers of 8 until the power is 1. The octal representation is formed by the quotients, written in the order generated by the algorithm. For example, to convert 12510 to octal: - 125 = 82 × 1 + 61 - 61 = 81 × 7 + 5 - 5 = 80 × 5 + 0 Therefore, 12510 = 1758. - 900 = 83 × 1 + 388 - 388 = 82 × 6 + 4 - 4 = 81 × 0 + 4 - 4 = 80 × 4 + 0 Therefore, 90010 = 16048. Method of successive multiplication by 8 To convert a decimal fraction to octal, multiply by 8; the integer part of the result is the first digit of the octal fraction. Repeat the process with the fractional part of the result, until it is null or within acceptable error bounds. Example: Convert 0.1640625 to octal: - 0.1640625 × 8 = 1.3125 = 1 + 0.3125 - 0.3125 × 8 = 2.5 = 2 + 0.5 - 0.5 × 8 = 4.0 = 4 + 0 Therefore, 0.164062510 = 0.1248. These two methods can be combined to handle decimal numbers with both integer and fractional parts, using the first on the integer part and the second on the fractional part. Method of successive duplication To convert integer decimals to octal, prefix the number with "0.". Perform the following steps for as long as digits remain on the right side of the radix: Double the value to the left side of the radix, using octal rules, move the radix point one digit rightward, and then place the doubled value underneath the current value so that the radix points align. If the moved radix point crosses over a digit that is 8 or 9, convert it to 0 or 1 and add the carry to the next leftward digit of the current value. Add octally those digits to the left of the radix and simply drop down those digits to the right, without modification. 0.4 9 1 8 decimal value +0 --------- 4.9 1 8 +1 0 -------- 6 1.1 8 +1 4 2 -------- 7 5 3.8 +1 7 2 6 -------- 1 1 4 6 6. octal value Octal to decimal conversion To convert a number k to decimal, use the formula that defines its base-8 representation: In this formula, ai is an individual octal digit being converted, where i is the position of the digit (counting from 0 for the right-most digit). Example: Convert 7648 to decimal: - 7648 = 7 × 82 + 6 × 81 + 4 × 80 = 448 + 48 + 4 = 50010 For double-digit octal numbers this method amounts to multiplying the lead digit by 8 and adding the second digit to get the total. Example: 658 = 6 × 8 + 5 = 5310 Method of successive duplication To convert octals to decimals, prefix the number with "0.". Perform the following steps for as long as digits remain on the right side of the radix: Double the value to the left side of the radix, using decimal rules, move the radix point one digit rightward, and then place the doubled value underneath the current value so that the radix points align. Subtract decimally those digits to the left of the radix and simply drop down those digits to the right, without modification. 0.1 1 4 6 6 octal value -0 ----------- 1.1 4 6 6 - 2 ---------- 9.4 6 6 - 1 8 ---------- 7 6.6 6 - 1 5 2 ---------- 6 1 4.6 - 1 2 2 8 ---------- 4 9 1 8. decimal value Octal to binary conversion To convert octal to binary, replace each octal digit by its binary representation. Example: Convert 518 to binary: - 58 = 1012 - 18 = 0012 Therefore, 518 = 101 0012. Binary to octal conversion The process is the reverse of the previous algorithm. The binary digits are grouped by threes, starting from the least significant bit and proceeding to the left and to the right. Add leading zeroes (or trailing zeroes to the right of decimal point) to fill out the last group of three if necessary. Then replace each trio with the equivalent octal digit. For instance, convert binary 1010111100 to octal: 001 010 111 100 1 2 7 4 Therefore, 10101111002 = 12748. Convert binary 11100.01001 to octal: 011 100 . 010 010 3 4 . 2 2 Therefore, 11100.010012 = 34.228. Octal to hexadecimal conversion The conversion is made in two steps using binary as an intermediate base. Octal is converted to binary and then binary to hexadecimal, grouping digits by fours, which correspond each to a hexadecimal digit. For instance, convert octal 1057 to hexadecimal: - To binary: 1 0 5 7 001 000 101 111 - then to hexadecimal: 0010 0010 1111 2 2 F Therefore, 10578 = 22F16. Hexadecimal to octal conversion Hexadecimal to octal conversion proceeds by first converting the hexadecimal digits to 4-bit binary values, then regrouping the binary bits into 3-bit octal digits. For example, to convert 3FA516: - To binary: 3 F A 5 0011 1111 1010 0101 - then to octal: 0 011 111 110 100 101 0 3 7 6 4 5 Therefore, 3FA516 = 376458. Due to having only factors of two, many octal fractions have repeating digits, although these tend to be fairly simple: Prime factors of the base: 2, 5 Prime factors of one below the base: 3 Prime factors of one above the base: 11 Other Prime factors: 7 13 17 19 23 29 31 Prime factors of the base: 2 Prime factors of one below the base: 7 Prime factors of one above the base: 3 Other Prime factors: 5 13 15 21 23 27 35 37 of the denominator \|Positional representation\|\|Positional representation\|\|Prime factors of the denominator \|1/3\|\|3\|\|0.3333... = 0.3\|\|0.2525... = 0.25\|\|3\|\|1/3\| \|1/6\|\|2, 3\|\|0.16\|\|0.125\|\|2, 3\|\|1/6\| \|1/10\|\|2, 5\|\|0.1\|\|0.06314\|\|2, 5\|\|1/12\| \|1/12\|\|2, 3\|\|0.083\|\|0.052\|\|2, 3\|\|1/14\| \|1/14\|\|2, 7\|\|0.0714285\|\|0.04\|\|2, 7\|\|1/16\| \|1/15\|\|3, 5\|\|0.06\|\|0.0421\|\|3, 5\|\|1/17\| \|1/18\|\|2, 3\|\|0.05\|\|0.034\|\|2, 3\|\|1/22\| \|1/20\|\|2, 5\|\|0.05\|\|0.03146\|\|2, 5\|\|1/24\| \|1/21\|\|3, 7\|\|0.047619\|\|0.03\|\|3, 7\|\|1/25\| \|1/22\|\|2, 11\|\|0.045\|\|0.02721350564\|\|2, 13\|\|1/26\| \|1/24\|\|2, 3\|\|0.0416\|\|0.025\|\|2, 3\|\|1/30\| \|1/26\|\|2, 13\|\|0.0384615\|\|0.02354\|\|2, 15\|\|1/32\| \|1/28\|\|2, 7\|\|0.03571428\|\|0.02\|\|2, 7\|\|1/34\| \|1/30\|\|2, 3, 5\|\|0.03\|\|0.02104\|\|2, 3, 5\|\|1/36\| The table below gives the expansions of some common irrational numbers in decimal and octal. \|√ (the length of the diagonal of a unit square)\|\|1.414213562373095048...\|\|1.3240 4746 3177 1674...\| \|√ (the length of the diagonal of a unit cube)\|\|1.732050807568877293...\|\|1.5666 3656 4130 2312...\| \|√ (the length of the diagonal of a 1×2 rectangle)\|\|2.236067977499789696...\|\|2.1706 7363 3457 7224...\| \|φ (phi, the golden ratio = (1+√)/2)\|\|1.618033988749894848...\|\|1.4743 3571 5627 7512...\| \|π (pi, the ratio of circumference to diameter of a circle)\|\|3.141592653589793238462643 \|3.1103 7552 4210 2643...\| \|e (the base of the natural logarithm)\|\|2.718281828459045235...\|\|2.5576 0521 3050 5355...\| - Computer numbering formats - Octal games, a game numbering system used in combinatorial game theory - Split octal, a 16-bit octal notation used by the Heath Company, DEC and others - Squawk code, a 12-bit octal representation of Gillham code - Syllabic octal, an octal representation of 8-bit syllables used by English Electric - Avelino, Heriberto (2006). "The typology of Pame number systems and the limits of Mesoamerica as a linguistic area" (PDF). Linguistic Typology. 10 (1): 41–60. doi:10.1515/LINGTY.2006.002. - Ascher, Marcia (1992). "Ethnomathematics: A Multicultural View of Mathematical Ideas". The College Mathematics Journal. 23 (4): 353–355. doi:10.2307/2686959. JSTOR 2686959. - Winter, Werner (1991). "Some thoughts about Indo-European numerals". In Gvozdanović, Jadranka (ed.). Indo-European numerals. Trends in Linguistics. 57. Berlin: Mouton de Gruyter. pp. 13–14. ISBN 3-11-011322-8. Retrieved 2013-06-09. - Wilkins, John (1668). An Essay Towards a Real Character and a Philosophical Language. London. p. 190. Retrieved 2015-02-08. - Donald Knuth, The Art of Computer Programming - See H. R. Phalen, "Hugh Jones and Octave Computation," The American Mathematical Monthly 56 (August–September 1949): 461-465. - James Anderson, On Octal Arithmetic [title appears only in page headers], Recreations in Agriculture, Natural-History, Arts, and Miscellaneous Literature, Vol. IV, No. 6 (February 1801), T. Bensley, London; pages 437-448. - Alfred B. Taylor, Report on Weights and Measures, Pharmaceutical Association, 8th Annual Session, Boston, 1859-09-15. See pages 48 and 53. - Alfred B. Taylor, Octonary numeration and its application to a system of weights and measures, Proc. Amer. Phil. Soc. Vol XXIV, Philadelphia, 1887; pages 296-366. See pages 327 and 330. - "TB Code Sheet". Dr. Dobb's Journal of Computer Calisthenics & Orthodontia, Running Light Without Overbyte. 1 (1). December 1975. - Microsoft Corporation (1987). "Constants, Variables, Expressions and Operators". GW-BASIC User's Manual. Retrieved 2015-12-12. - "2.4.1 Numeric Constants". CP/M-86 - Operating System - Programmer's Guide (PDF) (3 ed.). Pacific Grove, California, USA: Digital Research. January 1983 . p. 9. Archived (PDF) from the original on 2020-02-27. Retrieved 2020-02-27. (1+viii+122+2 pages) - Küveler, Gerd; Schwoch, Dietrich (2013) . Arbeitsbuch Informatik - eine praxisorientierte Einführung in die Datenverarbeitung mit Projektaufgabe (in German). Vieweg-Verlag, reprint: Springer-Verlag. doi:10.1007/978-3-322-92907-5. ISBN 978-3-528-04952-2. 978-3-32292907-5. Retrieved 2015-08-05. - Küveler, Gerd; Schwoch, Dietrich (2007-10-04). Informatik für Ingenieure und Naturwissenschaftler: PC- und Mikrocomputertechnik, Rechnernetze (in German). 2 (5 ed.). Vieweg, reprint: Springer-Verlag. ISBN 978-3-83489191-4. 978-3-83489191-4. Retrieved 2015-08-05. - Paul, Matthias R. (1997-07-30). NWDOS-TIPs — Tips & Tricks rund um Novell DOS 7, mit Blick auf undokumentierte Details, Bugs und Workarounds. MPDOSTIP. Release 157 (in German) (3 ed.). Archived from the original on 2016-11-04. Retrieved 2014-08-06. (NB. NWDOSTIP.TXT is a comprehensive work on Novell DOS 7 and OpenDOS 7.01, including the description of many undocumented features and internals. It is part of the author's yet larger MPDOSTIP.ZIPcollection maintained up to 2001 and distributed on many sites at the time. The provided link points to a HTML-converted older version of the - Paul, Matthias R. (2002-03-26). "Updated CLS posted". freedos-dev mailing list. Archived from the original on 2019-04-27. Retrieved 2014-08-06. - CCI Multiuser DOS 7.22 GOLD Online Documentation. Concurrent Controls, Inc. (CCI). 1997-02-10. HELP.HLP. - Haskell 98 Lexical Structure - OCaml: 7.1 Lexical conventions - Python 3: https://docs.python.org/3.1/reference/lexical_analysis.html#integer-literals - Perl 6: http://perlcabal.org/syn/S02.html#Radix_markers Archived 31 October 2014 at the Wayback Machine - RubySpec: https://github.com/kostya/rubyspec/blob/master/core/string/to_i_spec.rb[permanent dead link] - Tcl: http://wiki.tcl.tk/498 - ECMAScript 6th Edition draft: https://people.mozilla.org/~jorendorff/es6-draft.html#sec-literals-numeric-literals Archived 16 December 2013 at the Wayback Machine - "parseInt()", Mozilla Developer Network (MDN), If the input string begins with "0" (a zero), radix is assumed to be 8 (octal) or 10 (decimal). Exactly which radix is chosen is implementation-dependent. ECMAScript 5 clarifies that 10 (decimal) should be used, but not all browsers support this yet
18th and 19th centuriesThe town of Nashville was founded by James Robertson (explorer), James Robertson, John Donelson, and a party of Overmountain Men in 1779, near the original Cumberland settlement of Fort Nashborough. It was named for Francis Nash, the hero. Nashville quickly grew because of its strategic location as a port on the Cumberland River, a tributary of the Ohio River; and its later status as a major railroad center. By 1800, the city had 345 residents, including 136 enslaved African Americans and 14 free African Americans. In 1806, Nashville was Municipal corporation, incorporated as a city and became the of Davidson County, Tennessee. In 1843, the city was named as the permanent capital of the state of . The city government of Nashville owned 24 slaves by 1831, and 60 prior to the Civil War. They were "put to work to build the first successful water system and maintain the streets." Auction blocks and brokers' offices were part of the slave market at the heart of the city. It was the center of plantations cultivating tobacco and hemp as commodity crops, in addition to the breeding and training of thoroughbred horses, and other livestock. For years Nashville was considered one of the wealthiest southern capitals and a large portion of its prominence was from the iron business. Nashville led the south for iron production. The Nashville, Tennessee cholera epidemic (1849–50), cholera epidemic that struck Nashville in 1849–1850 took the life of former U.S. President James K. Polk and resulted in high fatalities. There were 311 deaths from cholera in 1849 and an estimated 316 to about 500 in 1850. Before the Civil War, about 700 free Blacks lived in small enclaves in northern Nashville while there were over 3,200 Black slaves in the city. By 1860, when the first American Civil War, rumblings of secession began to be heard across the Southern United States, South, antebellum Nashville was a prosperous city. The city's significance as a shipping port and rail center made it a desirable prize for competing military forces that wanted to control the region's important river and railroad transportation routes. In February 1862, Nashville became the first Confederate state capital to fall to Union (American Civil War), Union troops, and the state was occupied by Union troops for the duration of the war. Then African Americans from Middle Tennessee fled to contraband camps around military installations in Nashville's eastern, western, and southern borders. The Battle of Nashville (December 15–16, 1864) was a significant Union victory and perhaps the most decisive tactical victory gained by either side in the war; it was also the war's final major military action in which Tennessee regiments played a large part on both sides of the battle. Afterward, the Confederates conducted a war of attrition, making guerrilla raids and engaging in small skirmishes, with the Confederate forces in the Deep South almost constantly in retreat. In 1868, three years after the end of the Civil War, the Nashville chapter of the Ku Klux Klan was founded by Confederate veteran John W. Morton (Tennessee politician), John W. Morton. He was reported to have initiated General Nathan Bedford Forrest into the vigilante organization. Chapters of this secret insurgent group formed throughout the state and the South; they opposed voting and political organizing by freedmen, tried to control their behavior, and sometimes also attacked their White allies, including schoolteachers from the North. Whites directed violence against freedmen and their descendants both during and after the Reconstruction era. Two freedmen, Lynching of David Jones, David Jones and Lynching of Jo Reed, Jo Reed, were lynched in Nashville by White mobs in 1872 and 1875, respectively. Reed was hanged from a bridge over the river, but survived after the rope broke and he subsequently fell into the water, and fled the city soon thereafter. In 1873 Nashville suffered another cholera epidemic, along with towns throughout Sumner County along railroad routes and the Cumberland River. This was part of a larger epidemic that struck much of the United States. The epidemic is estimated to have killed around 1,000 people in Nashville. Meanwhile, the city had reclaimed its important shipping and trading position and developed a solid manufacturing base. The post–Civil War years of the late 19th century brought new prosperity to Nashville and Davidson County. Wealthy planters and businessmen built grand, classical-style buildings. A replica of the Parthenon (Nashville), Parthenon was constructed in Centennial Park (Nashville), Centennial Park, near downtown. On April 30, 1892, Ephraim Grizzard, an African American man, was lynching in the United States, lynched in a spectacle murder in front of a White mob of 10,000 in Nashville. His lynching was described by journalist Ida B. Wells as: "A naked, bloody example of the blood-thirstiness of the nineteenth century civilization of the Athens of the South." His brother, Henry Grizzard, had been lynched and hanged on April 24, 1892, in nearby Goodlettsville as a suspect in the same assault incident. From 1877 to 1950, a total of six lynchings of Blacks were conducted in Davidson County, four before the turn of the century. 20th centuryBy the turn of the century, Nashville had become the cradle of the Lost Cause of the Confederacy. The first chapter of the United Daughters of the Confederacy was founded here, and the ''Confederate Veteran'' magazine was published here. Most "guardians of the Lost Cause" lived Downtown or in the West End, near Centennial Park (Nashville), Centennial Park. At the same time, Jefferson Street (Nashville), Jefferson Street became the historic center of the African American community, with similar districts developing in the Black neighborhoods in East and North Nashville. In 1912, the Tennessee State University, Tennessee Agricultural and Industrial and Normal School as moved to Jefferson Street. The first Prince's Hot Chicken Shack originated at the corner of Jefferson Street and 28th Avenue in 1945. Jefferson Street became a destination for jazz and blues musicians, and remained so until the federal government split the area by construction of Interstate 40 in the late 1960s. In 1950 the state legislature approved a new city charter that provided for the election of city council members from s, rather than voting. This change was supported because at-large voting required candidates to gain a majority of votes from across the city. The previous system prevented the minority population, which then tended to support Republican candidates, from being represented by candidates of their choice; apportionment under single-member districts meant that some districts in Nashville had Black majorities. In 1951, after passage of the new charter, African American attorneys Z. Alexander Looby and Robert E. Lillard were elected to the city council. With the United States Supreme Court ruling in 1954 that public schools had to desegregate with "all deliberate speed", the family of student Robert Kelley filed a lawsuit in 1956, arguing that Nashville administrators should open all-White East High School to him. A similar case was filed by Reverend Henry Maxwell due to his children having to take a 45-minute bus ride from South Nashville to the north end of the city. These suits caused the courts to announce what became known as the "Nashville Plan", where the city's public schools would desegregate one grade per year beginning in the fall of 1957. Urban redevelopment accelerated over the next several decades, and the city grew increasingly segregated. An interstate was placed on the edge of East Nashville while another highway was built through Edgehill, a lower-income, predominantly minority community. Postwar development to presentRapid suburbanization occurred during the years immediately after World War II, as new housing was being built outside city limits. This resulted in a demand for many new schools and other support facilities, which the county found difficult to provide. At the same time, suburbanization led to a declining tax base in the city, although many suburban residents used unique city amenities and services that were supported financially only by city taxpayers. After years of discussion, a referendum was held in 1958 on the issue of consolidating city and county government. It failed to gain approval although it was supported by the then-elected leaders of both jurisdictions, County Judge Beverly Briley and Mayor Ben West. Following the referendum's failure, Nashville annexed some 42 square miles of suburban jurisdictions to expand its tax base. This increased uncertainty among residents, and created resentment among many suburban communities. Under the second charter for metropolitan government, which was approved in 1962, two levels of service provision were proposed: the General Services District and the Urban Services District, to provide for a differential in tax levels. Residents of the Urban Services District had a full range of city services. The areas that made up the General Services District, however, had a lower tax rate until full services were provided. This helped reconcile aspects of services and taxation among the differing jurisdictions within the large metro region. In the early 1960s, Tennessee still had racial segregation of public facilities, including lunch counters and department store fitting rooms. Hotels and restaurants were also segregated. Between February 13 and May 10, 1960, Nashville sit-ins, a series of sit-ins were organized at lunch counters in downtown Nashville by the Nashville Student Movement and Nashville Christian Leadership Council, and were part of a broader sit-in movement in the southeastern United States as part of an effort to end racial segregation of public facilities. On April 19, 1960, the house of Z. Alexander Looby, an African American attorney and council member, was bombed by segregationists. Protesters marched to the city hall the next day. Mayor Ben West said he supported the desegregation of lunch counters, which civil rights activists had called for. In 1963, Nashville consolidated its government with Davidson County, forming a Consolidated city-county, metropolitan government. The membership on the Metro Council, the legislative body, was increased from 21 to 40 seats. Of these, five members are elected and 35 are elected from s, each to serve a term of four years. In 1957 Nashville desegregated its school system using an innovative grade a year plan, in response to a class action suit Kelly vs. Board of Education of Nashville. By 1966 the Metro Council abandoned the grade a year plan and completely desegregated the entire school system at one time. Congress passed civil rights legislation in 1964 and 1965, but tensions continued as society was slow to change. On April 8, 1967, a riot broke out on the college campuses of Fisk University and Tennessee State University, historically black colleges and universities, historically Black colleges, after Stokely Carmichael spoke about Black Power at Vanderbilt University. Although it was viewed as a "race riot", it had classist characteristics. In 1979, the Ku Klux Klan burnt crosses outside two African American sites in Nashville, including the city headquarters of the NAACP. Since the 1970s the city and county have undergone tremendous growth, particularly during the boom and bust, economic boom of the 1990s under the leadership of then-Mayor and later-List of governors of Tennessee, Tennessee Governor, Phil Bredesen. Making urban renewal a priority, Bredesen fostered the construction or renovation of several city landmarks, including the Country Music Hall of Fame and Museum, the downtown Nashville Public Library, the Bridgestone Arena, and Nissan Stadium. Nissan Stadium (formerly Adelphia Coliseum and LP Field) was built after the National Football League's (NFL) Houston Oilers agreed to move to the city in 1995. The NFL team debuted in Nashville in 1998 at Vanderbilt Stadium, and Nissan Stadium opened in the summer of 1999. The Oilers changed their name to the Tennessee Titans and finished the season with the Music City Miracle and a close Super Bowl XXXIV, Super Bowl game. The St. Louis Rams won in the Final play of Super Bowl XXXIV, last play of the game. In 1997, Nashville was awarded a National Hockey League expansion team; this was named the Nashville Predators. Since the 2003–04 season, the Predators have made the playoffs in all but three seasons. In 2017, they made the Stanley Cup Finals for the first time in franchise history, but ultimately fell to the Pittsburgh Penguins, 4games to 2, in the best-of-seven series. 21st centuryOn January 22, 2009, residents rejected Nashville Charter Amendment 1, which sought to make English the official language of the city. Between May 1 and 7, 2010, much of Nashville was 2010 Tennessee floods, extensively flooded as part of a series of 1,000 year floods throughout Middle and West Tennessee. Much of the flooding took place in areas along the Cumberland and Harpeth River, Harpeth Rivers and Mill Creek (Davidson County, Tennessee), Mill Creek, and caused extensive damage to the many buildings and structures in the city, including the Grand Ole Opry House, Gaylord Opryland Resort & Convention Center, Opry Mills Mall, Schermerhorn Symphony Center, Bridgestone Arena, and Nissan Stadium. Sections of Interstate 24 and Briley Parkway were also flooded. Eleven people died in the Nashville area as a result of the flooding, and damages were estimated to be over $2 billion. The city recovered after the Great Recession. In March 2012, a Gallup poll ranked Nashville in the top five regions for job growth. In 2013, Nashville was described as "Nowville" and "It City" by ''GQ'', ''Forbes'', and ''The New York Times''. Nashville elected its first female mayor, Megan Barry, on September 25, 2015. As a council member, Barry had officiated at the city's first same-sex wedding on June 26, 2015. In 2017, Nashville's economy was deemed the third fastest-growing in the nation, and the city was named the "hottest housing market in the US" by Freddie Mac realtors. In May 2017, census estimates showed Nashville had passed Memphis, Tennessee, Memphis to become most populated city in Tennessee. Nashville has also made national headlines for its "homelessness crisis". Rising housing prices and the opioid crisis have resulted in more people being out on the streets: , between 2,300 and 20,000 Nashvillians are homeless. On March 6, 2018, due to felony charges filed against Mayor Barry relating to the misuse of public funds, she resigned before the end of her term. A 2018 Nashville mayoral special election, special election was called. Following a ruling by the Tennessee Supreme Court, the Davidson County Election Commission set the special election for May 24, 2018, to meet the requirement of 75 to 80 days from the date of resignation. David Briley, who was Vice Mayor during the Barry administration and Acting Mayor after her resignation, won the special election with just over 54% of the vote, becoming the 70th mayor of Nashville. On May 1, 2018, voters rejected Let's Move Nashville, a referendum which would have funded construction of an $8.9 billion mass transit system under the Nashville Metropolitan Transit Authority, by a 2 to 1 margin. On March 3, 2020, March 2020 Tennessee tornado outbreak, a tornado tracked west to east, just north of the downtown Nashville area, killing at least 25 people and leaving tens of thousands without electricity. Neighborhoods impacted included North Nashville, Germantown, East Nashville, Donelson, and Hermitage. On December 25, 2020, a 2020 Nashville bombing, vehicle exploded on Second Avenue, injuring three people. TopographyNashville lies on the Cumberland River in the northwestern portion of the Nashville Basin. Nashville's elevation ranges from its lowest point, Above mean sea level, above sea level at the Cumberland River, to its highest point, above sea level in the Radnor Lake State Natural Area. Nashville also sits at the start of the Highland Rim, a geophysical region of very hilly land. Because of this, Nashville is very hilly. Nashville also has some stand alone hills around the city such as the hill on which the Tennessee State Capitol building sits. According to the United States Census Bureau, the city has a total area of , of which of it is land and of it (4.53%) is water. CityscapeNashville's downtown area features a diverse assortment of entertainment, dining, cultural and architectural attractions. The Broadway and 2nd Avenue areas feature entertainment venues, night clubs and an assortment of restaurants. North of Broadway lie Nashville's central business district, Legislative Plaza, Capitol Hill and the Bicentennial Capitol Mall State Park, Tennessee Bicentennial Mall. Cultural and architectural attractions can be found throughout the city. Three major interstate highways (I-40, I-65 and I-24) converge near the core area of downtown, and many regional cities are within a day's driving distance. Nashville's first skyscraper, the Life & Casualty Tower, was completed in 1957 and launched the construction of other high rises in downtown Nashville. After the construction of the AT&T Building (Nashville), AT&T Building (commonly referred to by locals as the "Batman Building") in 1994, the downtown area saw little construction until the mid-2000s. The Pinnacle at Symphony Place, The Pinnacle, a high rise office building, opened in 2010, the first Nashville skyscraper completed in more than 15 years. Ten more skyscrapers have since been constructed or are under construction. Many civic and infrastructure projects are being planned, in progress, or recently completed. A new MTA bus hub was recently completed in downtown Nashville, as was the Music City Star pilot project. Several public parks have been constructed, such as the Public Square. Riverfront Park is scheduled to be extensively updated. The Music City Center opened in May 2013. It is a 1,200,000 square foot (110,000 m2) convention center with 370,000 square feet (34,000 m2) of exhibit space. FloraThe nearby city of Lebanon, Tennessee, Lebanon is notable and even named for its so-called "cedar glades", which occur on soils too poor to support most trees and are instead dominated by Juniperus virginiana, Virginian juniper. Blackberry bushes, Pinus virginiana, Virginia pine, Pinus taeda, loblolly pine, Sassafras albidum, sassafras, Acer rubrum, red maple, Betula nigra, river birch, Fagus grandifolia, American beech, Arundinaria gigantea, river cane, Kalmia latifolia, mountain laurel and Platanus occidentalis, sycamore are all common native trees, along with many others. In addition to the native forests, the combination of hot summers, abundant rainfall and mild winters permit a wide variety of both temperate and subtropical plants to be cultivated easily. Magnolia grandiflora, Southern magnolia and cherry blossom trees are commonly cultivated here, with the city having an annual cherry blossom festival. Lagerstroemia, Crepe myrtles and Taxus, yew bushes are also commonly grown throughout Metro Nashville, and the winters are mild enough that Magnolia virginiana, sweetbay magnolia is evergreen whenever it is cultivated. The pansy flower is popular to plant during the autumn, and some varieties will flower overwinter in Nashville's subtropical climate. However, many hot-weather plants like petunia and even Cyperus papyrus, papyrus thrive as annuals, and Musa basjoo, Japanese banana will die aboveground during winter but re-sprout after the danger of frost is over. Unbeknownst to most Tennesseans, even cold-hardy palms, particularly Rhapidophyllum, needle palm and Sabal minor, dwarf palmetto, are grown uncommonly but often successfully. High desert plants like Picea pungens, Colorado spruce and Opuntia humifusa, prickly pear cactus are also grown somewhat commonly, as are ''Yucca filamentosa''. ClimateNashville has a humid subtropical climate (Köppen climate classification, Köppen ''Cfa'', Trewartha climate classification, Trewartha ''Cf''), with hot, humid summers and generally cool winters typical of the Upper South. Monthly averages range from in January to in July, with a diurnal temperature variation of . Snowfall occurs during the winter months, but it is usually not heavy. Average annual snowfall is about , falling mostly in January and February and occasionally in March and December. The largest snow event since 2003 was on January 22, 2016, when Nashville received of snow in a single storm; the largest overall was , received on March 17, 1892, during the St. Patrick's Day Snowstorm. Rainfall is typically greater in November and December, and spring, while August to October are the driest months on average. Spring and fall are prone to severe thunderstorms, which may bring tornadoes, large hail, and damaging wind, with recent major events on Tornado outbreak of April 15–16, 1998, April 16, 1998; Tornado outbreak of April 6–8, 2006, April 7, 2006; 2008 Super Tuesday tornado outbreak, February 5, 2008; Tornado outbreak of April 9–11, 2009, April 10, 2009; 2010 Tennessee floods, May 1–2, 2010; and March 2020 Tennessee tornado outbreak, March 3, 2020. Relative humidity in Nashville averages 83% in the mornings and 60% in the afternoons, which is considered moderate for the Southeastern United States. In recent decades, due to urban development, Nashville has developed an urban heat island; especially on cool, clear nights, temperatures are up to warmer in the heart of the city than in rural outlying areas. The Nashville region lies within USDA Plant Hardiness Zone 7a. Nashville's long springs and autumns combined with a diverse array of trees and grasses can often make it uncomfortable for allergy sufferers. In 2008, Nashville was ranked as the 18th-worst spring allergy city in the U.S. by the Asthma and Allergy Foundation of America. The coldest temperature ever recorded in Nashville was on Winter 1985 cold wave, January 21, 1985, and the highest was on 2012 North American heat wave, June 29, 2012. DemographicsAccording to the 2016 American Community Survey, there were 667,885 people living in the city; in 2019 it rose to an estimated 670,820. The population density was . There were 294,794 housing units at an average density of . At the 2010 census, the racial makeup of the city was 65.5% White Americans, White (58.6% non-Hispanic White), 28.6% African American, 0.8% Native Americans in the United States, American Indian and Alaska Natives, Alaska Native, 3.5% Asian Americans, Asian, 0.1% Native Hawaiians, Native Hawaiian and Other Pacific Islander, and 1.4% from two or more races. 9.0% of the total population was of Hispanic and Latino Americans, Hispanic or Latino origin (they may be of any race). The non-Hispanic White population was 79.5% in 1970. There were 254,651 households and 141,469 families (55.6% of households). Of households with families, 37.2% had married couples living together, 14.1% had a female householder with no husband present, and 4.2% had a male householder with no wife present. 27.9% of all households had children under the age of 18, and 18.8% had at least one member 65 years of age or older. Of the 44.4% of households that are non-families, 36.2% were individuals, and 8.2% had someone living alone who was 65 years of age or older. The average household size was 2.38 and the average family size was 3.16. The age distribution was 22.2% under 18, 10.3% from 18 to 24, 32.8% from 25 to 44, 23.9% from 45 to 64, and 10.7% who were 65 or older. The median age was 34.2 years. For every 100 females, there were 94.1 males. For every 100 females age 18 and over, there were 91.7 males. The median income for a household in the city was $46,141, and the median income for a family was $56,377. Males with a year-round, full-time job had a median income of $41,017 versus $36,292 for females. The per capita income for the city was $27,372. About 13.9% of families and 18.2% of the population were below the poverty line, including 29.5% of those under age 18 and 9.9% of those age 65 or over. Of residents 25 or older, 33.4% have a bachelor's degree or higher. Because of its relatively low cost of living and large job market, Nashville has become a popular city for Immigration to the United States, immigrants. Nashville's foreign-born population more than tripled in size between 1990 and 2000, increasing from 12,662 to 39,596. The city's largest immigrant groups include Mexican Americans, Mexicans, Kurdish Americans, Kurds, Vietnamese Americans, Vietnamese, Laotian Americans, Laotians, Arab Americans, Arabs, and Somalis. There are also smaller communities of Pashtuns from Afghanistan and Pakistan concentrated primarily in Antioch, Tennessee, Antioch. Nashville has the largest History of the Kurds in Nashville, Tennessee, Kurdish community in the United States, numbering approximately 15,000. In 2009, about 60,000 Bhutanese refugees were being admitted to the U.S., and some were expected to resettle in Nashville. During the Iraqi parliamentary election, January 2005, Iraqi election of 2005, Nashville was one of the few international locations where Iraqi Americans, Iraqi expatriates could vote. The American Jews, American Jewish community in Nashville dates back over 150 years, and numbered about 8,000 in 2015, plus 2,000 Jewish college students. Metropolitan area, Nashville has the largest metropolitan area in the state of Tennessee, with an estimated population of 1,959,495. The Nashville metropolitan area encompasses 14 of 41 Middle Tennessee counties: Cannon County, Tennessee, Cannon, Cheatham County, Tennessee, Cheatham, Davidson County, Tennessee, Davidson, Dickson County, Tennessee, Dickson, Hickman County, Tennessee, Hickman, Macon County, Tennessee, Macon, Maury County, Tennessee, Maury, Robertson County, Tennessee, Robertson, Rutherford County, Tennessee, Rutherford, Smith County, Tennessee, Smith, Sumner County, Tennessee, Sumner, Trousdale County, Tennessee, Trousdale, Williamson County, Tennessee, Williamson, and Wilson County, Tennessee, Wilson. The 2019 population of the Nashville-Davidson–Murfreesboro–Columbia combined statistical area was estimated at 2,087,725. Religion59.6% of people in Nashville claim religious affiliation according to information compiled by Sperling's BestPlaces. The dominant religion in Nashville is Christianity, comprising 57.7% of the population. The Christian population is broken down into 20.6% Baptists, 6.2% Catholic Church, Catholics, 5.6% Methodism, Methodists, 3.4% Pentecostalism, Pentecostals, 3.4% Presbyterianism, Presbyterians, 0.8% Mormons, and 0.5% Lutheranism, Lutherans. 15.7% identify with other forms of Christianity, including the Eastern Orthodox Church, Orthodox Church and Christian Church (Disciples of Christ), Disciples of Christ. Islam is the second largest religion, comprising 0.8% of the population. 0.6% of the population adhere to eastern religions such as Buddhism, Sikhism, Jainism and Hinduism, and 0.3% follow Judaism. EconomyAs the "home of country music", Nashville has become a major music recording and production center. The music industry, Big Three record labels, as well as numerous independent labels, have offices in Nashville, mostly in the Music Row area. Nashville has been the headquarters of guitar company Gibson Guitar Corporation, Gibson since 1984. Since the 1960s, Nashville has been the second-largest music production center (after New York City) in the United States. Nashville's music industry is estimated to have a total economic impact of about $10billion per year and to contribute approximately 56,000 jobs to the Nashville area. In recent times Nashville has been described as a "southern boomtown" by numerous publications, with it having the third fastest growing economy in the United States as of 2017. It has been stated by the US Census bureau that Nashville "adds an average of 100 people a day to its net population increase". The Nashville region was also stated to be the "Number One" Metro Area for Professional and Business Service Jobs in America, as well as having the "hottest Housing market in America" as stated by Zillow. Although Nashville is renowned as a music recording center and tourist destination, its largest industry is health care. Nashville is home to more than 300 health care companies, including Hospital Corporation of America (HCA), the world's largest private operator of hospitals. , it is estimated the health care industry contributes per year and 200,000 jobs to the Nashville-area economy. CoreCivic, formerly known as Corrections Corporation of America and one of the largest private prison, private corrections company in the United States, was founded in Nashville in 1983, but relocated out of the city in 2019. Vanderbilt University was one of its investors prior to the company's initial public offering. The City of Nashville's pension fund includes "a $921,000 stake" in the company as of 2017. The ''Nashville Scene'' notes that, "A drop in CoreCivic stock value, however minor, would have a direct impact on the pension fund that represents nearly 25,000 current and former Metro employees." The automotive industry is also becoming important for the Middle Tennessee region. Nissan, Nissan North America moved its corporate headquarters in 2006 from Gardena, California (Los Angeles County, California, Los Angeles County) to Franklin, Tennessee, Franklin, a suburb south of Nashville. Nissan also has its largest North American manufacturing plant in Smyrna, Tennessee, Smyrna, another suburb of Nashville. Largely as a result of the increased development of Nissan and other Japanese economic interests in the region, Japan moved its former New Orleans consulate-general to Nashville's Palmer Plaza. General Motors also operates an Spring Hill Manufacturing, assembly plant in Spring Hill, Tennessee, Spring Hill, about south of Nashville. Bridgestone has a strong presence with their North American headquarters located in Nashville, with manufacturing plants and a distribution center in nearby counties. Other major industries in Nashville include insurance, finance, and publishing (especially religious publishing). The city hosts headquarters operations for several Protestant denominations, including the United Methodist Church, Southern Baptist Convention, National Baptist Convention, USA, Inc., National Baptist Convention USA, and the National Association of Free Will Baptists. Nashville is also known for some of their famously popular Southern confections, including Goo Goo Clusters (which have been made in Nashville since 1912). Fortune 500 companies with offices within Nashville include Bank of New York Mellon, BNY Mellon, Bridgestone Americas, Ernst & Young, Community Health Systems, Dell, Deloitte, Dollar General, Hospital Corporation of America, Nissan North America, Philips, Tractor Supply Company, and UBS. Of these, Community Health Systems, Dollar General, SmileDirectClub, Hospital Corporation of America, and Tractor Supply Company are headquartered in the city. In 2013, the city ranked No. 5 on ''Forbes'' list of the Best Places for Business and Careers. In 2015, Forbes put Nashville as the 4th Best City for White Collar Jobs. In 2015, Business Facilities' 11th Annual Rankings report named Nashville the number one city for Economic Growth Potential. In May 2018, AllianceBernstein pledged to build a private client office in the city by mid-2019 and to move its headquarters from New York City to Nashville by 2024. Additionally, in November 2018, Amazon (company), Amazon announced its plans to build an operations center in the Nashville Yards development to serve as the hub for their Retail Operations division. In December 2019, iHeartMedia selected Nashville as the site of its second digital headquarters. Real estate is becoming a driver for the city's economy. Based on a survey of nearly 1,500 real estate industry professionals conducted by PricewaterhouseCoopers and the Urban Land Institute, Nashville ranked 7th nationally in terms of attractiveness to real estate investors for 2016. , according to city figures, there is more than $2 billion in real estate projects underway or projected to start in 2016. Due to high yields available to investors, Nashville has been attracting a lot of capital from out-of-state. A key factor that has been attributed to the increase in investment is the adjustment to the city's zoning code. Developers can easily include a combination of residential, office, retail and entertainment space into their projects. Additionally, the city has invested heavily into public parks. Centennial Park is undergoing extensive renovations. The change in the zoning code and the investment in public space is consistent with the millennial generation's preference for walkable urban neighborhoods. Top employersAccording to the city's 2016 Comprehensive Annual Financial Report, the top employers in the city are: CultureMuch of the city's cultural life has revolved around its large university community. Particularly significant in this respect were two groups of critics and writers who were associated with Vanderbilt University in the early 20th century: the Fugitives (poets), Fugitives and the Southern Agrarians, Agrarians. Popular destinations include Fort Nashborough and Fort Negley, the former being a reconstruction of the original settlement, the latter being a semi-restored Civil War battle fort; the Tennessee State Museum; and Parthenon (Nashville), The Parthenon, a full-scale replica of the original Parthenon in Athens. The Tennessee State Capitol is one of the oldest working state capitol buildings in the nation. The Hermitage (Nashville, Tennessee), The Hermitage, the former home of President Andrew Jackson, is one of the largest presidential homes open to the public, and is also one of the most visited. DiningSome of the more popular types of local cuisine include hot chicken, hot fish, barbecue, and meat and three. Entertainment and performing artsNashville has a vibrant music and entertainment scene spanning a variety of genres. With a long history in the music scene it is no surprise that city was nicknamed 'Music City.' The Tennessee Performing Arts Center is the major performing arts center of the city. It is the home of the Nashville Repertory Theatre, the Nashville Opera Association, Nashville Opera, the Music City Drum and Bugle Corps, and the Nashville Ballet. In September 2006, the Schermerhorn Symphony Center opened as the home of the Nashville Symphony. As the city's name itself is a metonymy, metonym for the country music industry, many popular attractions involve country music, including the Country Music Hall of Fame and Museum, Belcourt Theatre, and Ryman Auditorium. Hence, the city became known as America's 'Country Music Capital.' The Ryman was home to the ''Grand Ole Opry'' until 1974 when the show moved to the Grand Ole Opry House, east of downtown. The Opry plays there several times a week, except for an annual winter run at the Ryman. Many music clubs and honky-tonk bars are in downtown Nashville, particularly the area encompassing Lower Broadway, Second Avenue, and Printer's Alley, which is often referred to as "the District". Each June, the CMA Music Festival (formerly known as Fan Fair) brings thousands of country fans to the city. The Tennessee State Fair is also held annually in September. Nashville was once home of television shows such as ''Hee Haw'' and ''Pop! Goes the Country'', as well as The Nashville Network and later, RFD-TV. Country Music Television and Great American Country currently operate from Nashville. The city was also home to the Opryland USA, Opryland USA theme park, which operated from 1972 to 1997 before being closed by its owners (Gaylord Entertainment Company) and soon after demolished to make room for the Opry Mills mega-shopping mall. The Contemporary Christian music industry is based along Nashville's Music Row, with a great influence in neighboring Williamson County, Tennessee, Williamson County. The Christian record companies include EMI Christian Music Group, Provident Label Group and Word Records. Music Row houses many gospel music and Contemporary Christian music companies centered around 16th and 17th Avenues South. On River Road, off Charlotte Pike in West Nashville, the ''CabaRay'' opened its doors on January 18, 2018. The performing venue of Ray Stevens, it offers a Vegas-style dinner and a show atmosphere. There is also a piano bar and a gift shop. Although Nashville was never known as a major jazz town, it did have many great jazz bands, including The Nashville Jazz Machine led by Dave Converse and its current version, the Nashville Jazz Orchestra, led by Jim Williamson, as well as The Establishment, led by Billy Adair. The Francis Craig Orchestra entertained Nashvillians from 1929 to 1945 from the Oak Bar and Grille Room in the Hermitage Hotel. Craig's orchestra was also the first to broadcast over local radio station WSM-AM and enjoyed phenomenal success with a 12-year show on the NBC Red Network, NBC Radio Network. In the late 1930s, he introduced a newcomer, Dinah Shore, a local graduate of Hume Fogg High School and Vanderbilt University. Radio station WMOT, WMOT-FM in nearby Murfreesboro, Tennessee, Murfreesboro, which formerly programmed jazz, aided significantly in the recent revival of the city's jazz scene, as has the non-profit Nashville Jazz Workshop, which holds concerts and classes in a renovated building in the north Nashville neighborhood of Germantown. Fisk University also maintains a jazz station, WFSK. Nashville has an active theatre scene and is home to several professional and community theatre companies. Nashville Children's Theatre, Nashville Repertory Theatre, the Nashville Shakespeare Festival, the Dance Theatre of Tennessee and the Tennessee Women's Theater Project are among the most prominent professional companies. One community theatre, Circle Players, has been in operation for over 60 years. The Barbershop Harmony Society has its headquarters in Nashville. TourismPerhaps the biggest factor in drawing visitors to Nashville is its association with country music, in which the Nashville sound played a role. Many visitors to Nashville attend live performances of the Grand Ole Opry, the world's longest-running live radio show. The Country Music Hall of Fame and Museum is another major attraction relating to the popularity of country music. The Gaylord Opryland Resort & Convention Center, the Opry Mills regional shopping mall and the ''General Jackson (riverboat), General Jackson'' showboat, are all located in what is known as Music Valley. Civil War history is important to the city's tourism industry. Sites pertaining to the Battle of Nashville and the nearby Battle of Franklin (1864), Battle of Franklin and Battle of Stones River can be seen, along with several well-preserved antebellum plantation houses such as Belle Meade Plantation, Carnton plantation in Franklin, and Belmont Mansion. Nashville has many arts centers and museums, including the Frist Center for the Visual Arts, Cheekwood Botanical Garden and Museum of Art, the Tennessee State Museum, the Johnny Cash Museum, Fisk University's Van Vechten and Aaron Douglas Galleries, Vanderbilt University's Fine Art Gallery and Sarratt Gallery, the National Museum of African American Music, and the Parthenon (Nashville), full-scale replica of the Parthenon. Nashville has become an increasingly popular destination for bachelor party, bachelor and bachelorette party, bachelorette parties. In 2017 ''Nashville Scene'' counted 33 bachelorette parties on Lower Broadway ("from Fifth Avenue down to the Cumberland River, it's their town") in less than two hours on a Friday night, and stated that the actual number was likely higher. Downtown, the newspaper wrote, "offers five blocks of bars with live music and no cover". In 2018, ''The New York Times'' called Nashville "the hottest destination for bachelorette parties in the country" because of the honky-tonk bars' live music. City boosters welcome the bachelorette parties because temporary visitors may become permanent; ''BuzzFeed'' wrote, "These women are at precisely the point in their lives when a move to Nashville is possible". The CMT (U.S. TV channel), CMT reality television series ''Bachelorette Weekend'' follows the employees at Bach Weekend, a Nashville company that designs and throws bachelor and bachelorette parties. Major annual events NicknamesNashville is a colorful, well-known city in several different arenas. As such, it has earned various sobriquets, including: * Music City, U.S.A.: WSM (AM), WSM-AM announcer David Cobb first used this name during a 1950 broadcast and it stuck. It is now the official nickname used by the Nashville Convention and Visitors Bureau. Nashville is the home of the Grand Ole Opry, the Country Music Hall of Fame, and many major record labels. This name also dates back to 1873, where after receiving and hearing a performance by the Fisk Jubilee Singers, Queen Victoria of the United Kingdom is reported as saying that "These young people must surely come from a musical city." Athens of the South: Home to 24 post-secondary educational institutions, Nashville has long been compared to Athens, the ancient city of learning and site of Plato's Academy. Since 1897, a full-scale Parthenon (Nashville), replica of the Athenian Parthenon has stood in Nashville, and many examples of classical and neoclassical architecture can be found in the city. The term was popularized by Philip Lindsley (1786–1855), President of the University of Nashville, though it is unclear whether he was the first person to use the phrase. The Protestant Vatican City, Vatican or The Bible Belt#Buckle, Buckle of the Bible Belt: Nashville has over 700 churches, several seminaries, a number of Christian music companies, and is the headquarters for the publishing arms of the Southern Baptist Convention (LifeWay Christian Resources), the United Methodist Church (United Methodist Publishing House) and the National Baptist Convention, USA, Inc., National Baptist Convention (Sunday School Publishing Board). It is also the seat of the National Baptist Convention, the National Association of Free Will Baptists, the Gideons International, the Gospel Music Association, and Thomas Nelson (publisher), Thomas Nelson, the world's largest producer of Bibles. * Cashville: Nashville native Young Buck released a successful rap album called ''Straight Outta Cashville'' that has popularized the nickname among a new generation. * Little Kurdistan: Nashville has the United States' largest population of History of the Kurds in Nashville, Tennessee, Kurdish people, estimated to be around 11,000. * Nash Vegas or Nashvegas Nashville has additionally earned the moniker "The Hot Chicken Capital", becoming known for the local specialty cuisine hot chicken. The Music City Hot Chicken Festival is hosted annually in Nashville and several restaurants make this spicy version of southern fried chicken. ProfessionalNashville is home to four professional sports franchises. Three play at the major professional sports leagues in the United States and Canada, highest professional level of their respective sports: the Tennessee Titans of the National Football League (NFL), the Nashville Predators of the National Hockey League (NHL), and Nashville SC of Major League Soccer (MLS). The city is also home to one minor league team: the Nashville Sounds of Minor League Baseball's Triple-A East. An investment group, Music City Baseball, seeks to secure a Major League Baseball expansion franchise or lure an existing team to the city. The Tennessee Titans moved to Nashville in 1998. Previously known as the History of the Houston Oilers, Houston Oilers, which began play in 1960 in Houston, Houston, Texas, the team relocated to Tennessee in 1997. They played at the Liberty Bowl Memorial Stadium in Memphis, Tennessee, Memphis for one season, then moved to Nashville in 1998 and played in Vanderbilt Stadium for one season. During those two years, the team was known as the Tennessee Oilers, but changed its name to Titans in 1999. The team now plays at Nissan Stadium in Nashville, which opened in 1999. Since moving to Nashville, the Titans have won four division championships (2000, 2002, 2008, and 2020) and one conference championship (1999). They competed in 1999's Super Bowl XXXIV, losing to the History of the St. Louis Rams, St. Louis Rams, 23–16. The city previously hosted the 1939 Nashville Rebels of the American Football League (1938), American Football League and two Arena Football League teams named the Nashville Kats (1997–2001 and 2005–2007). From April 25–27, 2019, Nashville hosted the 2019 NFL Draft, which saw an estimated 200,000 fans attend each day. The Nashville Predators joined the National Hockey League as an expansion team in the 1998–99 season. The team plays its home games at Bridgestone Arena. The Predators have won two division championships (2017–18 and 2018–19) and one conference championship (2016–17). Nashville SC, a Major League Soccer franchise, began play in 2020 at Nissan Stadium. It is expected to relocate to the Nashville Fairgrounds Stadium upon its planned completion in 2022. The Nashville Sounds baseball team was established in 1978 as expansion franchise of the Double-A (baseball), Double-A Southern League (baseball), Southern League. The Sounds won the league championship in 1979 and 1982. In 1985, the Double-A Sounds were replaced by a Triple-A (baseball), Triple-A team of the American Association (20th century), American Association. After the American Association dissolved in 1997, the Sounds joined the Triple-A Pacific Coast League in 1998 and won the league championship in 2005. The Sounds left their original ballpark, Herschel Greer Stadium, in 2015 for First Horizon Park, a new ballpark built on the site of the former Sulphur Dell ballpark. They became members of the Triple-A East in 2021. In total, the Sounds have won ten division titles and three league championships. Nashville is the home of the second-oldest continually operating race track in the United States, the Fairgrounds Speedway. It hosted NASCAR Winston Cup races from 1958 to 1984, NASCAR Busch Series and NASCAR Truck Series in the 1980s and 1990s, and later the NASCAR Whelen All-American Series and ARCA Racing Series. Nashville Superspeedway is located southeast of Nashville in Gladeville, Tennessee, Gladeville, part of the Nashville Metropolitan Statistical Area. The track held NASCAR sanctioned events from 2001 to 2011 as well as IndyCar races from 2001 to 2008. Nashville Superspeedway will reopen in 2021 and host the premier NASCAR Cup Series for the first time. The Nashville Invitational was a golf tournament on the PGA Tour from 1944 to 1946. The Sara Lee Classic was part of the LPGA Tour from 1988 to 2002. The BellSouth Senior Classic of the Champions Tour was held from 1994 to 2003. The Nashville Golf Open is part of the Web.com Tour since 2016. The 1961 Women's Western Open and 1980 U.S. Women's Open were also held in Nashville. College and amateurNashville is also home to four NCAA Division I, Division I athletic programs. Nashville is also home to the National Collegiate Athletic Association, NCAA college football Music City Bowl. The Nashville Rollergirls are Nashville's only women's flat track roller derby team. Established in 2006, Nashville Rollergirls compete on a regional and national level. They play their home games at the Nashville Fairgrounds Sports Arena. In 2014, they hosted the WFTDA Championships at Municipal Auditorium. The Nashville Kangaroos are an Australian Rules Football team that compete in the United States Australian Football League. The Kangaroos play their home games at Elmington Park. The team is the reigning USAFL Central Region Champions. Three Little League Baseball teams from Nashville (one in 1970 Little League World Series, 1970; one in 2013 Little League World Series, 2013; and, one in 2014 Little League World Series, 2014) have qualified for the Little League World Series. Teams from neighboring Goodlettsville qualified for the 2012 Little League World Series, 2012 and 2016 Little League World Series, 2016 series, giving the metropolitan area teams in three consecutive years to so qualify; and four teams in five years. Parks and gardensNashville Board of Parks and Recreation, Metro Board of Parks and Recreation owns and manages of land and 99 parks and greenways (comprising more than 3% of the total area of the county). Warner Parks, situated on of land, consists of a learning center, of scenic roads, of hiking trails, and of horse trails. It is also the home of the annual Iroquois (horse), Iroquois Steeplechase (horse racing), Steeplechase. The United States Army Corps of Engineers maintains parks on Old Hickory Lake and Percy Priest Lake. These parks are used for activities such as fishing, water skiing, sailing and boating. The Harbor Island Yacht Club makes its headquarters on Old Hickory Lake, and Percy Priest Lake is home to the Vanderbilt Sailing Club and Nashville Shores. Other parks in Nashville include Centennial Park (Nashville), Centennial Park, Shelby Park (Nashville), Shelby Park, Cumberland Park, and Radnor Lake State Natural Area. On August 27, 2013, Nashville mayor Karl Dean revealed plans for two new riverfront parks on the east and west banks of the Cumberland River downtown. Construction on the east bank park began in the fall of 2013, and the projected completion date for the west bank park is 2015. Among many exciting benefits of this Cumberland River re-development project is the construction of a highly anticipated outdoor amphitheater. Located on the west bank, this music venue will be surrounded by a new park and will replace the previous thermal plant site. It will include room for 6,500 spectators with 2,500 removable seats and additional seating on an overlooking grassy knoll. In addition, the east bank park will include a river landing, providing people access to the river. In regard to the parks' benefits for Nashvillian civilians, Mayor Dean remarked that "if done right, the thermal site can be an iconic park that generations of Nashvillians will be proud of and which they can enjoy". Law and governmentThe city of Nashville and Davidson County merged in 1963 as a way for Nashville to combat the problems of urban sprawl. The combined entity is officially known as "the Metropolitan Government of Nashville and Davidson County", and is popularly known as "Metro Nashville" or simply "Metro". It offers services such as Metropolitan Nashville Police Department, police, Nashville Fire Department, fire, Nashville Electric Service, electricity, water and sewage treatment. When the Metro government was formed in 1963, the government was split into two service districts—the "urban services district" and the "general services district." The urban services district encompasses the 1963 boundaries of the former City of Nashville, approximately , and the general services district includes the remainder of Davidson County. There are six smaller municipalities within the consolidated city-county: Belle Meade, Tennessee, Belle Meade, Berry Hill, Tennessee, Berry Hill, Forest Hills, Tennessee, Forest Hills, Oak Hill, Tennessee, Oak Hill, Goodlettsville, Tennessee, Goodlettsville (partially), and Ridgetop, Tennessee, Ridgetop (partially). These municipalities use a two-tier system of government, with the smaller municipality typically providing police services and the Metro Nashville government providing most other services. Previously, the city of Lakewood, Tennessee, Lakewood also had a separate charter. However, Lakewood residents voted in 2010 and 2011 to dissolve its city charter and join the metropolitan government, with both votes passing. Nashville is governed by a mayor, vice-mayor, and Metropolitan Council of Nashville and Davidson County, 40-member Metropolitan Council. It uses the strong-mayor form of the mayor–council government, mayor–council system. The current mayor of Nashville is John Cooper (Tennessee politician), John Cooper. The Metropolitan Council is the legislative body of government for Nashville and Davidson County. There are five council members who are elected at large and 35 council members that represent individual districts. The Metro Council has regular meetings that are presided over by the vice-mayor, who is currently Jim Shulman. The Metro Council meets on the first and third Tuesday of each month at 6:00pm, according to the Metropolitan Charter. Nashville is home to the Tennessee Supreme Court's courthouse for Middle Tennessee and the Estes Kefauver Federal Building and United States Courthouse, home of the United States District Court for the Middle District of Tennessee. PoliticsNashville has been a Democratic Party (United States), Democratic stronghold since at least the end of Reconstruction era, Reconstruction, and has remained staunchly Democratic even as the state as a whole has trended strongly Republican Party (United States), Republican. Pockets of Republican influence exist in the wealthier portions of the city, but they are usually no match for the overwhelming Democratic trend in the rest of the city. The issue of school busing roiled politics for years but subsided after the 1990s. While local elections are officially nonpartisan, nearly all the city's elected officials are publicly known as Democrats. The city is split among 10 state house districts, all of which are held by Democrats. Three state senate districts and part of a fourth are within the county; two are held by Democrats and two by Republicans. In the state legislature, Nashville politicians serve as leaders of both the Tennessee Senate, Senate and Tennessee House of Representatives, House Democratic Caucuses. Representative Mike Stewart serves as Chairman of the House Caucus. Senator Jeff Yarbro serves as Chairman of the Senate Caucus. Democrats are no less dominant at the federal level. Democratic presidential candidates have failed to carry Davidson County only five times since reconstruction; in 1928, 1968, 1972, 1984 and 1988.David Leip's Presidential Atlas (Maps for Indiana by election) CrimeAccording to the FBI's Uniform Crime Reporting database, Metropolitan Nashville has a violent crime rate approximately three times the national average, and a property crime rate approximately 1.6 times the average. The following table shows Nashville's crime rate per 100,000 inhabitants for seven UCR categories. EducationThe city is served by Metropolitan Nashville Public Schools, also referred to as Metro Schools. This district is the second largest school district in Tennessee, and enrolls approximately 85,000 students at 169 schools. In addition, Nashville is home to numerous private schools, including Montgomery Bell Academy, Harpeth Hall School, University School of Nashville, Lipscomb Academy, The Ensworth School, Christ Presbyterian Academy, Father Ryan High School, Pope John Paul II High School (Tennessee), Pope John Paul II High School, Franklin Road Academy, Davidson Academy, Nashville Christian School, Donelson Christian Academy, and St. Cecilia Academy (Nashville, Tennessee), St. Cecilia Academy. Combined, all of the private schools in Nashville enroll more than 15,000 students. Colleges and universitiesNashville is often labeled the "Athens of the South" due to the many colleges and universities in the city and the metropolitan area. Vanderbilt University is the largest university in Nashville, with approximately 13,000 students. Vanderbilt is considered one of the nation's leading research universities and is particularly known for its medical, law, and education programs. Nashville is home to four historically black colleges and universities, historically Black institutions of higher education, the second highest in the nation, behind Atlanta, Atlanta, Georgia. These are Fisk University, Tennessee State University, Meharry Medical College, and American Baptist College. Other schools based in Nashville include Belmont University, Lipscomb University, Trevecca Nazarene University, John A. Gupton College. The Tennessee Board of Regents operates Nashville State Community College and the Nashville branch of the Tennessee Colleges of Applied Technology. In total, enrollment in post-secondary education in Nashville is around 43,000. In addition, there are several other institutes of higher education in the Nashville metropolitan area. Middle Tennessee State University (MTSU), a full-sized public university with Tennessee's second largest undergraduate population, is located in Murfreesboro. Other schools include Daymar College and O'More College of Design, both in Franklin, and Cumberland University in Lebanon, Tennessee, Lebanon. MediaThe daily newspaper in Nashville is ''The Tennessean'', which until 1998 competed with the ''Nashville Banner'', another daily paper that was housed in the same building under a joint operating agreement, joint-operating agreement. ''The Tennessean'' is the city's most widely circulated newspaper. Online news service ''NashvillePost.com'' competes with the printed dailies to break local and state news. Several weekly papers are also published in Nashville, including ''The Nashville Pride'', ''Nashville Business Journal'', ''Nashville Scene'' and ''The Tennessee Tribune''. Historically, ''The Tennessean'' was associated with a broadly liberal editorial policy, while ''The Banner'' carried staunchly conservative views in its editorial pages; ''The Banner''s heritage had been carried on, to an extent, by ''The City Paper'' which folded in August 2013 after having been founded in October 2000. The ''Nashville Scene'' is the area's alternative newspaper, alternative weekly broadsheet. ''The Nashville Pride'' is aimed towards community development and serves Nashville's entrepreneurial population. ''Nashville Post'' is an online news source covering business, politics and sports. Nashville is home to eleven broadcast television stations, although most households are served by direct Cable television, cable network connections. Comcast, Comcast Cable has a monopoly on terrestrial cable service in Davidson County (but not throughout the entire media market). Nashville is ranked as the 29th largest television market in the United States. Major stations include WKRN-TV 2 (American Broadcasting Company, ABC), WSMV-TV 4 (NBC), WTVF 5 (CBS), WNPT (TV), WNPT 8 (PBS), WZTV 17 (Fox Broadcasting Company, Fox), WNPX-TV 28 (Ion Television, ion), WPGD-TV 50 (Trinity Broadcasting Network, TBN), WLLC-LP 42 (Univision), WUXP-TV 30 (MyNetworkTV), (WJFB) 44 (MeTV), and WNAB 58 (The CW, CW). Nashville is also home to cable networks Country Music Television (CMT), among others. CMT's master control facilities are located in New York City with the other Viacom (2005–present), Viacom properties. The Top 20 Countdown and CMT Insider are taped in their Nashville studios. Shop at Home Network was once based in Nashville, but the channel signed off in 2008. Several Frequency modulation, FM and Amplitude modulation, AM radio stations broadcast in the Nashville area, including five Campus radio, college stations and one low-power broadcasting, LPFM community radio station. Nashville is ranked as the 44th largest radio market in the United States. WSM-FM is owned by Cumulus Media and is 95.5 FM. WSM (AM), WSM-AM, owned by Gaylord Entertainment Company, can be heard nationally on 650 AM or online a TransportationAccording to the 2016 American Community Survey, 78.1% of working Nashville residents commuted by driving alone, 9.8% carpooled, 2% used public transportation, and 2.2% walked. About 1.1% used all other forms of transportation, including taxicab, motorcycle, and bicycle. About 6.7% of working Nashville residents worked at home. In 2015, 7.9% of city of Nashville households were without a car; this figure decreased to 5.9% in 2016. The national average was 8.7 percent in 2016. Nashville averaged 1.72 cars per household in 2016, compared to a national average of 1.8 per household. HighwaysNashville is centrally located at the crossroads of three Interstate Highway System, Interstate Highways, Interstate 40, I-40 (east-west), Interstate 24, I-24 (northwest-southeast) and Interstate 65, I-65 (north-south). I-40 connects the city between Memphis, Tennessee, Memphis to the west and Knoxville, Tennessee, Knoxville to the east, I-24 connects between Clarksville, Tennessee, Clarksville to the northwest and Chattanooga, Tennessee, Chattanooga to the southeast, and I-65 connects between Louisville, Kentucky to the north and Huntsville, Alabama to the south. All three of these interstate highways, which also serve the suburbs, form brief concurrency (road), concurrencies with each other in the city, and completely encircle downtown. Interstate 440 (Tennessee), Interstate 440 is a bypass route connecting I-40, I-65, and I-24 south of downtown Nashville. Tennessee State Route 155, Briley Parkway, the majority of which is a controlled access highway, freeway, forms a bypass around the north side of the city and its interstates. Ellington Parkway, a freeway made up of a section of U.S. Route 31E in Tennessee, U.S. Route 31E, runs between east of downtown and Briley Parkway, serving as an alternative route to I-65. Interstate 840 (Tennessee), Interstate 840 provides an outer southern bypass for the city and its suburbs. United States Numbered Highway System, U.S. Routes U.S. Route 31 in Tennessee, 31, 31E, U.S. Route 31W in Tennessee, 31W, U.S. Route 31A, 31 Alternate, U.S. Route 41 in Tennessee, 41, U.S. Route 41 Alternate (Monteagle, Tennessee–Hopkinsville, Kentucky), 41 Alternate, U.S. Route 70 in Tennessee, 70, U.S. Route 70S, 70S, and U.S. Route 431 in Tennessee, 431 also serve Nashville, intersecting in the city's center as arterial surface roads and radiating outward. Most of these routes are called "pikes" and many carry the names of nearby towns to which they lead. Among these are Clarksville Pike, Gallatin Pike, Lebanon Pike, Murfreesboro Pike, Nolensville Pike, and Franklin Pike. Public transitThe Nashville Metropolitan Transit Authority, Metropolitan Transit Authority provides bus transit within the city. Routes utilize a spoke–hub distribution paradigm, hub and spoke method, centered around the Music City Central transit station in downtown. A rejected Let's Move Nashville, expansion plan included use of bus rapid transit and light rail service at some point in the future. Nashville is considered a gateway city for rail and air traffic for the Piedmont Atlantic Megaregion. AirThe city is served by Nashville International Airport (BNA), which is operated by the Metropolitan Nashville Airport Authority (MNAA). 18.27 million passengers visited the airport in 2019, making it the 31st busiest airport in the US. BNA is ranked fastest growing airport among the top 50 airports in the United States. Nashville International Airport serves 600 daily flights to more than 85 nonstop markets. In late 2014, BNA became the first major U.S. airport to establish dedicated pick-up and drop-off areas for vehicle for hire companies. The airport authority also operates the John C. Tune Airport, a Class E airspace general aviation airport. Intercity railAlthough a major freight hub for CSX Transportation, Nashville is not currently served by Amtrak, the List of major cities in U.S. lacking inter-city rail service, third-largest metropolitan area in the U.S. to have this distinction. Nashville's Union Station (Nashville), Union Station had once been a major intercity passenger rail center for the Louisville and Nashville Railroad; Nashville, Chattanooga and St. Louis Railway; and the Tennessee Central Railway, reaching Midwestern cities and cities on the Gulf of Mexico and the Atlantic Ocean. However, by the time of Amtrak's founding, service had been cut back to a single train, the ''Floridian (train), Floridian'', which ran from Chicago to Miami and St. Petersburg, Florida. It served Union Station until its cancellation on October 9, 1979, due to poor track conditions resulting in late trains and low ridership, ending over 120 years of intercity rail service in Nashville. While there have been few proposals to restore Amtrak service to Nashville, there have been repeated calls from residents. In addition to scarce federal funding, Tennessee state officials do not believe that Nashville is large enough to support intercity rail. "It would be wonderful to say I can be in Memphis and jump on a train to Nashville, but the volume of people who would do that isn't anywhere close to what the cost would be to provide the service," said Ed Cole, chief of environment and planning with the Tennessee Department of Transportation. Ross Capon, executive director of the National Association of Railroad Passengers, said rail trips would catch on if routes were expanded, but conceded that it would be nearly impossible to resume Amtrak service to Nashville without a substantial investment from the state. However, in 2020, Amtrak indicated it was considering a service that would run from Atlanta to Nashville by way of Chattanooga. Nashville launched a passenger commuter rail system called the Music City Star on September 18, 2006. The only currently operational leg of the system connects the city of Lebanon to downtown Nashville at the Nashville Riverfront station. Legs to Clarksville, Murfreesboro and Gallatin are currently in the feasibility study stage. The system plan includes seven legs connecting Nashville to surrounding suburbs. BridgesBridges within the city include: UtilitiesThe city of Nashville owns the Nashville Electric Service (NES), Metro Water Services (MWS) and Nashville District Energy System (NDES). The Nashville Electric Service provides electricity to the entirety of Davidson County and small portions of the six adjacent counties, and purchases its power from the Tennessee Valley Authority. Metro Water Services provides water, wastewater, and stormwater to Nashville and the majority of Davidson County, as well as water services to small portions of Rutherford and Williamson counties, and wastewater services to small portions of all of the surrounding counties except for Cheatham County. MWS sources its water from the Cumberland River and operates two water treatment plants and three wastewater treatment plants. Ten additional utility companies provide water and sewer service to Nashville and Davidson County. The Nashville District Energy System provides heating and cooling services to certain buildings in downtown, including multiple government buildings. Natural gas is provided by Piedmont Natural Gas, a subsidiary of Duke Energy. HealthcareAs a major center for the healthcare industry, Nashville is home to several hospitals and other primary care facilities. Most hospitals in Nashville are operated by Vanderbilt University Medical Center, the TriStar Division of Hospital Corporation of America, and Saint Thomas Health. The Metropolitan Nashville Hospital Authority operates Nashville General Hospital, which is affiliated with Meharry Medical College. Sister citiesNashville's Sister city, sister cities are: * Belfast, Northern Ireland * Caen, France * Chengdu, China * Edmonton, Canada * Kamakura, Japan * Magdeburg, Germany * Mendoza, Argentina, Mendoza, Argentina * Taiyuan, China * Tamworth Regional Council, Tamworth, Australia ;Candidates * Gwangjin District, Gwangjin (Seoul), South Korea ;International Friendship City * Crouy, France See also* List of people from Nashville, Tennessee * ''The Children (Halberstam), The Children'', 1999 book about the Nashville Student Movement * National Register of Historic Places listings in Davidson County, Tennessee Further reading* * * * * * * * * * * * * * * * *
Telephony is the field of technology involving the development, application, and deployment of telecommunication services for the purpose of electronic transmission of voice, fax, or data, between distant parties. The history of telephony is intimately linked to the invention and development of the telephone. Telephony is commonly referred to as the construction or operation of telephones and telephonic systems and as a system of telecommunications in which telephonic equipment is employed in the transmission of speech or other sound between points, with or without the use of wires. The term is also used frequently to refer to computer hardware, software, and computer network systems, that perform functions traditionally performed by telephone equipment. In this context the technology is specifically referred to as Internet telephony, or voice over Internet Protocol (VoIP). The first telephones were connected directly in pairs. Each user had a separate telephone wired to the locations he might wish to reach. This quickly became inconvenient and unmanageable when people wanted to communicate with more than a few people. The inventions of the telephone exchange provided the solution for establishing telephone connections with any other telephone in service in the local area. Each telephone was connected to the exchange via one wire pair, the local loop. Nearby exchanges in other service areas were connected with trunk lines and long distance service could be established by relaying the calls through multiple exchanges. Initially the switchboards were manually operated by an attendant, a switchboard operator. When a customer cranked a handle on the telephone, it turned on an indicator on the board in front of the operator who would plug the operator headset into that jack and offer service. The caller had to ask for the called party by name, later by number, and the operator connected one end of a circuit into the called party jack to alert them. If the called station answered the operator disconnected their headset and complete the station-to-station circuit. Trunk calls were made with the assistance of other operators at other exchangers in the network. In modern times, most telephones are plugged into telephone jacks. The jacks are connected by inside wiring to a drop wire which connects the building to a cable. Cables usually bring a large number of drop wires from all over a district access network to one wire center or telephone exchange. When a telephone user wants to make a telephone call, equipment at the exchange examines the dialed telephone number and connects that telephone line to another in the same wire center, or to a trunk to a distant exchange. Most of the exchanges in the world are interconnected through a system of larger switching systems, forming the public switched telephone network (PSTN). After the middle of the 20th century, fax and data became important secondary users of the network created to carry voices, and late in the century, parts of the network were upgraded with ISDN and DSL to improve handling of such traffic. Today, telephony uses digital technology (digital telephony) in the provisioning of telephone services and systems. Telephone calls can be provided digitally, but may be restricted to cases in which the last mile is digital, or where the conversion between digital and analog signals takes place inside the telephone. This advancement has reduced costs in communication, and improved the quality of voice services. The first implementation of this, ISDN, permitted all data transport from end-to-end speedily over telephone lines. This service was later made much less important due to the ability to provide digital services based on the IP protocol. Since the advent of personal computer technology in the 1980s, computer telephony integration has progressively provided more sophisticated telephony services, initiated and controlled by the computer, such as making and receiving voice, fax, and data calls with telephone directory services and caller identification. The integration of telephony software and computer systems is a major development in the evolution of the automated office. The term is used in describing the computerized services of call centers, such as those that direct your phone call to the right department at a business you're calling. It's also sometimes used to describe the ability to use your personal computer to initiate and manage phone calls (in which case you can think of your computer as your personal call center). CTI is not a new concept and has been used in the past in large telephone networks, but only dedicated call centers could justify the costs of the required equipment installation. Primary telephone service providers are offering information services such as automatic number identification, which is a telephone service architecture that separates CTI services from call switching and will make it easier to add new services. Dialed Number Identification Service (DNIS) on a scale is wide enough for its implementation to bring real value to business or residential telephone usage. A new generation of applications (middleware) is being developed as a result of standardization and availability of low cost computer telephony links. The term's scope has been broadened with the advent of the different new communication technologies. In its broadest sense, the terms encompasses phone communication, Internet calling, mobile communication, faxing, voicemail and video conferencing. Telephony's initial idea returns to POTS, (an acronym for "plain old telephone service") technically called the PSTN (public-switched telephone network). This system is being fiercely challenged by and to a great extent yielding to Voice over IP (VoIP) technology, which is also commonly referred to as IP Telephony and Internet Telephony. IP telephony is a modern form of telephony which uses the TCP/IP protocol popularized by the Internet to transmit digitized voice data. Also, unlike traditional phone service, IP telephony service is relatively unregulated by government. In the United States, the Federal Communications Commission (FCC) regulates phone-to-phone connections, but says they do not plan to regulate connections between a phone user and an IP telephony service provider.Using the Internet, calls travel as packets of data on shared lines, avoiding the tolls of the PSTN. The challenge in IP telephony is to deliver the voice, fax, or video packets in a dependable flow to the user. Much of IP telephony focuses on that challenge. Starting with the introduction of the transistor, invented in 1947 by Bell Laboratories, to amplification and switching circuits in the 1950s, and through development of computer-based electronic switching systems, the public switched telephone network (PSTN) has gradually evolved towards automation and digitization of signaling and audio transmissions. Digital telephony is the use of digital electronics in the operation and provisioning of telephony systems and services. Since the 1960s a digital core network has replaced the traditional analog transmission and signaling systems, and much of the access network has also been digitized. Digital telephony has dramatically improved the capacity, quality, and cost of the network. End-to-end analog telephone networks were first modified in the early 1960s by upgrading transmission networks with Digital Signal 1 (DS1/T1) carrier systems, designed to support the basic 3 kHz voice channel by sampling the bandwidth-limited analog voice signal and encoding using PCM. While digitization allows wideband voice on the same channel, the improved quality of a wider analog voice channel did not find a large market in the PSTN. Later transmission methods such as SONET and fiber optic transmission further advanced digital transmission. Although analog carrier systems existed that multiplexed multiple analog voice channels onto a single transmission medium, digital transmission allowed lower cost and more channels multiplexed on the transmission medium. Today the end instrument often remains analog but the analog signals are typically converted to digital signals at the serving area interface (SAI), central office (CO), or other aggregation point. Digital loop carriers (DLC) place the digital network ever closer to the customer premises, relegating the analog local loop to legacy status. Milestones in digital telephony - early experiments with pulse code modulation in telephony - the 8-bit, 8 kHz standard is developed; Nyquist's theorem and the standard 3.5 kHz telephony bandwidth - DS0 as the basic digital telephony bitstream standard - non-linear quantization: A-law vs. ?-law, and transcoding between the two - bit error rate and intelligibility - first practical digital telephone systems put into service - the U.S. T-carrier system and the European E-carrier system developed to carry digital telephony - introduction of space-time switching in fully digital electronic switching systems - replacement of tone signaling with digital signaling for trunks - in-band signaling vs. out-of-band signaling - the problem of bit-robbing - development of SS7 - emergence of fiber optic networking allows greater reliability and call capacity - transition from plesiochronous transmission to synchronous systems like SONET/SDH - optical self-healing ring networks further increase reliability - digital/optical systems revolutionize international long-distance networks, particularly undersea cables - digital telephone exchanges eliminate moving parts, make exchange equipment much smaller and more reliable - separation of exchange and concentrator functions - roll-out of digital systems throughout the PSTN - provision of intelligent network services - digital speech coding and compression - speech compression on international digital trunks - phone tapping in the digital environment - introduction of digital mobile telephony, specialized compression algorithms for high bit error rates - direct digital termination to customers via ISDN; PRI catches on, BRI mostly does not, except in Germany - the effects of digital telephony, and digital termination at the ISP, on modem performance - voice over IP as a carrier strategy - emergence of ADSL leads to voice over IP becoming a consumer product, and the slow demise of dial-up Internet access - expected convergence of VoIP, mobile telephony, etc. - flattening of telephony tariffs, increasing moves towards flat rate pricing as the marginal cost of telephony drops further and further. Main article: Voice over IP A commercial IP telephone, with keypad, control keys, and screen functions to perform configuration and user features. A specialization of digital telephony, Internet Protocol (IP) telephony involves the application of digital networking technology that was the foundation to the Internet to create, transmit, and receive telecommunications sessions over computer networks. Internet telephony is commonly known as voice over Internet Protocol (VoIP), reflecting the principle, but it has been referred with many other terms. VoIP has proven to be a disruptive technology that is rapidly replacing traditional telephone infrastructure technologies. As of January 2005, up to 10% of telephone subscribers in Japan and South Korea have switched to this digital telephone service. A January 2005 Newsweek article suggested that Internet telephony may be "the next big thing". As of 2006, many VoIP companies offer service to consumers and businesses. IP telephony uses an Internet connection and hardware IP phones, analog telephone adapters, or softphone computer applications to transmit conversations encoded as data packets. In addition to replacing plain old telephone service (POTS), IP telephony services compete with mobile phone services by offering free or lower cost connections via WiFi hotspots. VoIP is also used on private networks which may or may not have a connection to the global telephone network.
The scientific importance of these first samples from the Galaxy can’t be overstated. The major form of heavy elements in interstellar space is in dust. This interstellar dust plays a major role in the formation of new stars and planetary systems. Our own Solar System formed from gas and dust in the interstellar medium 4.6 billion years ago. The heavy elements making up Earth and our bodies were once interstellar dust. In the words of Joni Mitchell, “We are Stardust.” But we don’t even know what the typical interstellar dust grain looks like! We are extremely excited about the prospect of studying directly the first contemporary interstellar dust. Interstellar dust was first discovered flowing across the Solar System by dust detectors aboard the Ulysses spacecraft in 1993 and was later confirmed by the Galileo mission to Jupiter. The particles were identified as coming from a location in the sky in the Constellation Ophiuchus, looking toward the center of the Milky Way Galaxy. What we have learned about interstellar dust comes from remote observations of how the dust absorbs, scatters, polarizes, and even emits light. Also, some ancient interstellar dust has been identified in meteorites found on Earth. Interstellar dust is small, ranging in size from 0.01 microns all the way up to 20 microns. They are made of different minerals such as silicates, graphitic carbon, hydrogenated amorphous carbon, alumina, and even diamond carbon. Interstellar dust grains form by condensation in the regions around stars that are coming to the end of their life cycle: red giants, planetary nebulae, white dwarfs, novae, and supernovae. The dust grains mix with the interstellar medium (the stuff between the stars) and slowly experience chemical and isotopic changes from interactions with the gas and radiation in interstellar space. Dust grains do not last for very long in the interstellar medium before being dissociated by radiation, maybe a few hundreds of millions of years. This is why we say that the dust collected by the Stardust mission is contemporary dust, it must be only a few hundred million years old at most, whereas dust found recovered in meteorites would have been incorporated into them at the time of the formation of the Solar System (4.6 billion years ago). While finding the interstellar dust grains captured in Stardusts aerogel collectors is the goal of the Stardust@home project, the identification of these grains is only the first step. The next step is the analysis. Once we have a few examples to examine, a committee of experts will decide on the next steps. Because they are so small and so precious, each track is worth about a million dollars if there turn out to be 100 of them! The analysis of these particles will have to be done extremely carefully and will take many years. Many types of analysis destroy the samples, so we will have to start with the gentlest techniques and proceed very carefully. The great advantage of this type of sample-return mission is that one can take advantage of the improvements in analytical techniques for years or even decades to come. Analytical techniques improved dramatically even during the seven years between the launch and the return of Stardust, and there is no sign of a slowdown in progress. So no matter what, some of these interstellar dust particles will be set aside for our great-grandchildren to analyze. More on interstellar dust from the JPL Stardust website: Aerogel is one of the strangest materials ever developed. It is a solid, yet is only a few times as dense as air. If you hold it in your hand, you can only barely feel its weight, and it looks bluish and ghostly like solid smoke. While it looks blue, it casts an orange shadow. It does this for the same reason that the sky is blue and sunsets are red! Aerogel has extremely bizarre properties. It is a solid, glassy nanofoam, yet weighs next to nothing. Aerogel has the almost magical property that it can capture particles moving at very high speeds (several miles per second or more) better than any other material. In some cases, particles can be captured in a nearly pristine state. Particles moving at these speeds vaporize if they hit any other material.
A Kater's pendulum is a reversible freeswinging pendulum invented by British physicist and army captain Henry Kater in 1817 for use as a gravimeter instrument to measure the local acceleration of gravity. Its advantage is that, unlike previous pendulum gravimeters, the pendulum's centre of gravity and center of oscillation do not have to be determined, allowing greater accuracy. For about a century, until the 1930s, Kater's pendulum and its various refinements remained the standard method for measuring the strength of the Earth's gravity during geodetic surveys. It is now used only for demonstrating pendulum principles. So by measuring the length L and period T of a pendulum, g can be calculated. The Kater pendulum consists of a rigid metal bar with two pivot points, one near each end of the bar. It can be suspended from either pivot and swung. It also has either an adjustable weight that can be moved up and down the bar, or one adjustable pivot, to adjust the periods of swing. In use, it is swung from one pivot, and the period timed, and then turned upside down and swung from the other pivot, and the period timed. The movable weight (or pivot) is adjusted until the two periods are equal. At this point the period T is equal to the period of an 'ideal' simple pendulum of length equal to the distance between the pivots. From the period and the measured distance L between the pivots, the acceleration of gravity can be calculated with great precision from the equation (1) above. Gravity measurement with pendulums The first person to discover that gravity varied over the Earth's surface was French scientist Jean Richer, who in 1671 was sent on an expedition to Cayenne, French Guiana, by the French Académie des Sciences, assigned the task of making measurements with a pendulum clock. Through the observations he made in the following year, Richer determined that the clock was 2½ minutes per day slower than at Paris, or equivalently the length of a pendulum with a swing of one second there was 1¼ Paris lines, or 2.6 mm, shorter than at Paris. It was realized by the scientists of the day, and proven by Isaac Newton in 1687, that this was due to the fact that the Earth was not a perfect sphere but slightly oblate; it was thicker at the equator because of the Earth's rotation. Since the surface was farther from the Earth's center at Cayenne than at Paris, gravity was weaker there. Since that time pendulums began to be used as precision gravimeters, taken on voyages to different parts of the world to measure the local gravitational acceleration. The accumulation of geographical gravity data resulted in more and more accurate models of the overall shape of the Earth. Pendulums were so universally used to measure gravity that, in Kater's time, the local strength of gravity was usually expressed not by the value of the acceleration g now used, but by the length at that location of the seconds pendulum, a pendulum with a period of two seconds, so each swing takes one second. It can be seen from equation (1) that for a seconds pendulum, the length is simply proportional to g: Inaccuracy of gravimeter pendulums In Kater's time, the period T of pendulums could be measured very precisely by timing them with precision clocks set by the passage of stars overhead. Prior to Kater's discovery, the accuracy of g measurements was limited by the difficulty of measuring the other factor L, the length of the pendulum, accurately. L in equation (1) above was the length of an ideal mathematical 'simple pendulum' consisting of a point mass swinging on the end of a massless cord. However the 'length' of a real pendulum, a swinging rigid body, known in mechanics as a compound pendulum, is more difficult to define. In 1673 Dutch scientist Christiaan Huygens in his mathematical analysis of pendulums, Horologium Oscillatorium, showed that a real pendulum had the same period as a simple pendulum with a length equal to the distance between the pivot point and a point called the center of oscillation, which is located under the pendulum's center of gravity and depends on the mass distribution along the length of the pendulum. The problem was there was no way to find the location of the center of oscillation in a real pendulum accurately. It could theoretically be calculated from the shape of the pendulum if the metal parts had uniform density, but the metallurgical quality and mathematical abilities of the time didn't allow the calculation to be made accurately. To get around this problem, most early gravity researchers, such as Jean Picard (1669), Charles Marie de la Condamine (1735), and Jean-Charles de Borda (1792) approximated a simple pendulum by using a metal sphere suspended by a light wire. If the wire had negligible mass, the center of oscillation was close to the center of gravity of the sphere. But even finding the center of gravity of the sphere accurately was difficult. In addition, this type of pendulum inherently wasn't very accurate. The sphere and wire didn't swing back and forth as a rigid unit, because the sphere acquired a slight angular momentum during each swing. Also the wire stretched elastically during the pendulum's swing, changing L slightly during the cycle. However, in Horologium Oscillatorium, Huygens had also proved that the pivot point and the center of oscillation were interchangeable. That is, if any pendulum is suspended upside down from its center of oscillation, it has the same period of swing, and the new center of oscillation is the old pivot point. The distance between these two conjugate points was equal to the length of a simple pendulum with the same period. As part of a committee appointed by the Royal Society in 1816 to reform British measures, Kater had been contracted by the House of Commons to determine accurately the length of the seconds pendulum in London. He realized Huygens principle could be used to find the center of oscillation, and so the length L, of a rigid (compound) pendulum. If a pendulum were hung upside down from a second pivot point that could be adjusted up and down on the pendulum's rod, and the second pivot were adjusted until the pendulum had the same period as it did when swinging right side up from the first pivot, the second pivot would be at the center of oscillation, and the distance between the two pivot points would be L. Kater wasn't the first to have this idea. French mathematician Gaspard de Prony first proposed a reversible pendulum in 1800, but his work was not published till 1889. In 1811 Friedrich Bohnenberger again discovered it, but Kater independently invented it and was first to put it in practice. Kater built a pendulum consisting of a brass rod about 2 meters long, 1½ inches wide and one-eighth inch thick, with a weight (d) on one end. For a low friction pivot he used a pair of short triangular 'knife' blades attached to the rod. In use the pendulum was hung from a bracket on the wall, supported by the edges of the knife blades resting on flat agate plates. The pendulum had two of these knife blade pivots (a), facing one another, about a meter (40 in) apart, so that a swing of the pendulum took approximately one second when hung from each pivot. Kater found that making one of the pivots adjustable caused inaccuracies, making it hard to keep the axis of both pivots precisely parallel. Instead he permanently attached the knife blades to the rod, and adjusted the periods of the pendulum by a small movable weight (b,c) on the pendulum shaft. Since gravity only varies by a maximum of 0.5% over the Earth, and in most locations much less than that, the weight only had to be adjusted slightly. Moving the weight toward one of the pivots decreased the period when hung from that pivot, and increased the period when hung from the other pivot. This also had the advantage that the precision measurement of the separation between the pivots only had to be made once. To use, the pendulum was hung from a bracket on a wall, with the knife blade pivots supported on two small horizontal agate plates, in front of a precision pendulum clock to time the period. It was swung first from one pivot, and the oscillations timed, then turned upside down and swung from the other pivot, and the oscillations timed again. The small weight (c) was adjusted with the adjusting screw, and the process repeated until the pendulum had the same period when swung from each pivot. By putting the measured period T, and the measured distance between the pivot blades L, into the period equation (1), g could be calculated very accurately. Kater performed 12 trials. He measured the period of his pendulum very accurately using the clock pendulum by the method of coincidences; timing the interval between the coincidences when the two pendulums were swinging in synchronism. He measured the distance between the pivot blades with a microscope comparator, to an accuracy of 10−4 in. (2.5 μm). As with other pendulum gravity measurements, he had to apply small corrections to the result for a number of variable factors: - the finite width of the pendulum's swing, which increased the period - temperature, which caused the length of the rod to vary due to thermal expansion - atmospheric pressure, which reduced the effective mass of the pendulum by the buoyancy of the displaced air, increasing the period - altitude, which reduced the gravitational force with distance from the center of the Earth. Gravity measurements are always referenced to sea level. He gave his result as the length of the seconds pendulum. After corrections, he found that the mean length of the solar seconds pendulum at London, at sea level, at 62 °F (17 °C), swinging in vacuum, was 39.1386 inches. This is equivalent to a gravitational acceleration of 9.81158 m/s2. The largest variation of his results from the mean was 0.00028 inches (7.1 µm). This represented a precision of gravity measurement of 7(10−6) (7 milligals). In 1824, the British Parliament made Kater's measurement of the seconds pendulum the official standard of length for defining the yard. The large increase in gravity measurement accuracy made possible by Kater's pendulum established gravimetry as a regular part of geodesy. To be useful, it was necessary to find the exact location (latitude and longitude) of the 'station' where a gravity measurement was taken, so pendulum measurements became part of surveying. Kater's pendulums were taken on the great historic geodetic surveys of much of the world that were being done during the 19th century. In particular, Kater's pendulums were used in the Great Trigonometric Survey of India. Repeatedly timing each period of a Kater pendulum, and adjusting the weights until they were equal, was time consuming and error-prone. Friedrich Bessel showed in 1826 that this was unnecessary. As long as the periods measured from each pivot, T1 and T2, are close in value, the period T of the equivalent simple pendulum can be calculated from them: Here and are the distances of the two pivots from the pendulum's center of gravity. The distance between the pivots, , can be measured with great accuracy. and , and thus their difference , cannot be measured with comparable accuracy. They are found by balancing the pendulum on a knife edge to find its center of gravity, and measuring the distances of each of the pivots from the center of gravity. However, because is so much smaller than , the second term on the right in the above equation is small compared to the first, so doesn't have to be determined with high accuracy, and the balancing procedure described above is sufficient to give accurate results. Therefore, the pendulum doesn't have to be adjustable at all, it can simply be a rod with two pivots. As long as each pivot is close to the center of oscillation of the other, so the two periods are close, the period T of the equivalent simple pendulum can be calculated with equation (2), and the gravity can be calculated from T and L with (1). In addition, Bessel showed that if the pendulum was made with a symmetrical shape, but internally weighted on one end, the error caused by effects of air resistance would cancel out. Also, another error caused by the finite diameter of the pivot knife edges could be made to cancel out by interchanging the knife edges. Bessel didn't construct such a pendulum, but in 1864 Adolf Repsold, under contract to the Swiss Geodetic Commission, developed a symmetric pendulum 56 cm long with interchangeable pivot blades, with a period of about ¾ second. The Repsold pendulum was used extensively by the Swiss and Russian Geodetic agencies, and in the Survey of India. Other widely used pendulums of this design were made by Charles Peirce and C. Defforges. - Kater, Henry (1818). "An account of experiments for determining the length of the pendulum vibrating seconds in the latitude of London". Phil. Trans. R. Soc. London. 104 (33): 109. Retrieved 2008-11-25. - Nave, C. R. (2005). "Simple Pendulum". Hyperphysics. Dept. of Physics and Astronomy, Georgia State Univ. Retrieved 2009-02-20. - Poynting, John Henry; Joseph John Thompson (1907). A Textbook of Physics, 4th Ed. London: Charles Griffin & Co. p. 20. - Victor F., Lenzen; Robert P. Multauf (1964). "Paper 44: Development of gravity pendulums in the 19th century". United States National Museum Bulletin 240: Contributions from the Museum of History and Technology reprinted in Bulletin of the Smithsonian Institution. Washington: Smithsonian Institution Press. p. 307. Retrieved 2009-01-28. - Zupko, Ronald Edward (1990). Revolution in Measurement: Western European Weights and Measures since the Age of Science. New York: Diane Publishing. pp. 107–110. ISBN 0-87169-186-8. - Lenzen & Multauf 1964, p. 315 - Poynting & Thompson 1907, p. 12 - Elias Loomis (1864). Elements of Natural Philosophy, 4th Ed. New York: Harper & Brothers. p. 109. - Torge, Wolfgang (2001). Geodesy: An Introduction. Walter de Gruyter. p. 177. ISBN 3-11-017072-8. - Poynting & Thompson 1907, p. 15 - The Accurate Measurement of g using Kater's pendulum, U. of Sheffield Has derivation of equations - Kater, Henry (June 1818) An Account of the Experiments for determining the length of the pendulum vibrating seconds in the latitude of London, The Edinburgh Review, Vol. 30, p.407 Has detailed account of experiment, description of pendulum, value determined, interest of French scientists
It was the Japanese carrier attack on Pearl Harbor that brought America into World War II. Had the Japanese not attacked, it is unclear just when America would have entered the War. The Japanese Imperial Fleet was a superbly trained force with modern, well designed vessels. Many naval experts at the time did not fully appreciate the effectivness of the Imperial Navy. The lack of radar, however, proved a huge disadvantage. Allied radar and many other technical advances were the result of close cooperation between American and British scientists anf joint development projects that began even before America entered the War. There was no comparable Axis technical cooperation or even coordination of military campaigns. The Kriegsmarine had very effective radar on its surface ships like Bismarck yet advanced German technology like radar, jet engines, and other equipment was not provided to the Japanese until very late in the War, too late to be of any effective use to the Japanese war effort. While Pearl Harbor was a stunning tactical victory, it was a strategic blunder by the Japanese of incaluable proportions. The Japanese were able to seize much of Southeast Asia, but the stunning American carrier victory at Midway, significantly reduced the strike capability of the Imperial Navy. This provided the time for American industrial capacity to reated a naval force with which Japan's limited industrial capacity could not cope. While the German submarine campaign in the North Atlantic failed, the American submarine campaign in thePacific proved spectacularly successful. The Japanese merchant marine was almost completely destroying, cutting the country's war industries off from supplies and bringing the country close to starvation. Amercan industrial strength enabled America to build a naval force capable of leap froging from island to island. The Navy by 1944 had seized islands from which the Japanese Home Island could be bombed. The Navy also enabled the Army to retake New Guinea and the Phillipines and by 1945 Okinawa. Navy and Army forces were preparing for a full-scale amphibious invasion of the Home Islands when two atomic bombs were dropped (August 1945) and Japan finally surrendered (September 1945). The Pacific is by far the world's largest ocean. The expanse and the distances involved are beyond the experience of most people. Only people who have flown from the States to Australia and New Zealand begin to get some idea of the distances involved. The great expanse is only broken by tiny islands and atolls, around which is where the major battles of the Pacific War were fought. The Pacific is both the largest and deepest of the five oceans. It extends from the Arctic Ocean south to the Antarctic. It is ringed by Asia and Australia in the west and North and South America in the east. The separation of Asia and the Americas is immense. There is a reason that Columbus not only made it home, but crossed the Atlantic several times. But Magellan did not even make it all the way across the Pacific. The Pacific is also the world's deepest ocean but that was not a factor in the Pacific War--the overriding factor was the immense area involved. Ironically, the Japanese attack on Pearl Harbor suceeded because The U.S. Navy believed that the Pacific distances made an attack impossible. And the Japanese believed that the distances made an attack on their country impossible. Even a child can marvel at the lack of insight to think that Japan could attack the Americans, but the Americans could not attack Japan. Actually, at the time of the Japanese Pearl Harbor attack, work on the American long-range B-29 Superfortress was already underway. The Pacific exceeds 165 million square kilometers. Numbers alone are difficult to comprehend. The Pcific encompases an eaexceeding that of the earth's entire land area. This is about a third of the world's surface and about half of the water surface and total quantity of free water. The oceanic pole of inaccessibility is in the Pacific Ocean. The Battle of the Atlantic was fought over a small area beyond air cover -- the Mid-Atlantic Gap or Black Pit. In the Pacific, it was not a small area of the Pacific beyond air cover -- it was virtually the entire expanse. This is why the Pacific War became an island hopping campaigns beginning with Guadalcanal in the Solomons--the furthest point of Japanese expansion. . The islands were needed to expand air cover. Here the Japanese had an advantage because of the longer range of their aircraft. It was an advantage, however, that did not provide them the ability to deliver a knock-out blow and prevent America from mobilizing for war. There are several marginal seas in the Pacific, including the Bismarck Sea, Coral Sea, East China Sea, Java Sea, Philippine Sea, Sea of Japan, Sea of Okhotsk, Solomon Sea, South China Sea, and the Tasman Sea. This includes the area of the Western Pacific where most of the Pacific War battles were fought. The only exceptions are when the Imperial Fleet ventured east to attack the Americans. Given the huge expanse it is stunning that when American carriers suddenly appeared at just the precise time and place (Coral Sea and Midway) to engage Imperial Navy thrusts, the Japanese did not immediately understand that their codes had been broken. And even more mystifying after studing the matter, dismissed the possibility. The Americans capitalized on the huge distances by waging a submarine campaign on Japanese sea lanes, destroying the maru fleet. Japan had a sizeable submarine fleet, but ignored the German advice and declined to attack the American sea lane China for several centurues after the first European ships arrived restricted foreign trders desiring to do business in China. One major concern was that the Europeans had little the Chinese wanted. This ended with the British Opium Wars (mid-19th century), forcing the Chinese to permit imports of opium. The British and other Europeans forced the Chinese to open their ports. The Europeans also forced territorial concessions and extra territoriality. The Chinese who had for centuries been the dominant power in Asia, suddenly found that their armies and navy were impotent against European forces with modern weapons. The Chinese lost the Sino-Japanese War (1894-95) with the rapidly modernizing Japanese. The Boxer Rebellion was a reflection of Chinese frustration (1900), but only depened the decline of Chinese power. The Europeans controlled Hong Kong, Wei Hai Wei, and Tsingtau, or had concessions in ports like Canton and Shanghai. The United States promoted the Open Door Policy, but was concerned with its interests. The European powers to protect their commercial interests and citizens in China maintained naval forces in various places in the Far East, including China. This became known as China Station. The United States had an interest in trade as well as after the Spanish American War (1898-9) possession of the Philippines also maintained a squadron in China. Of particular concern was Shanghai, China's principal port, located at the mouth of the Yangtze River which led into the inteior. To ensure that the river was kept open, a new clas of vessel with shallow drafts were developed--China gunboats. The purpose was to 'show the flag', fight pirates, and protect foreign-owned vessels plying the river. Small ships like destroyers and mindsweepers could also enter the lower Yangtze up to Nanking. The China Station proved to be colorful duty, but there were dangers. The Yangtze is one of the world's geat rivers. It is also wild and unpredictable in many areas. Small craft like gunboats could be driven ashore or smashed on the rocks in the towering gorges. Pirates were a continuing problem in the unsettled conditions of the early-20th century. And this increased as the Imperial regime came apart. Then there was for several years the problem of war lords. There were several incidents during the 190s when foreign ships were attacked. There were several such attacks during 1927. Japan invaded China (1937). Foreigners were caught between Chinese and Japanese forces, especially in the first year of the War which involved conventional operations and centered on Shanghai and the Yagstze river ports leading to Nanking. The Americans and Europeans continued to maintain a presence in China as best they could. The Panay incident was a clear indication of Japanese intentions (1937). European warships called at Shanghai and were fired on or threatened by the Japanese. With the outbreak of World War Ii in Europe, the British and French graduallywithdrew removed their forces from China. Only leaving only a few small vessels to show the flag. As the European and American presence waned and the Japanese occupied most of coastal China, they became increasingly assertive, imposing control over movement of foreign warships in Chinese waters. The United States by late-1941 also largely withdrew. The major powers after World War I chastened by the incredible loss of life and destruction persued a policy of naval disarmament. The Treaties flowing from the Washington Naval Conference (1921-22) limited national fleets. Interestingly, Admiral Yamamoto suggested that battleships be scsrapped. I am not sure what the Japanese onjective was with this proposal. It was not taken seriously by the other naval powers. One result was that the Japanese began to take an interest in carrier which were not covered by the treaties. A naval building program was persued in viloation of the treaties, although I am not sure when this begun. Certianly the construction of Mustashi? and Yamato far exceeded the limitations (1937). The Japanese Navy was disturbed that the Japanese did not receive parity with the American and British fleet. Many advocated an aggressive preparation for war. Others such as Admiral Isoroku Yamamoto insisted that war with America and Britian would be suicidal because of their superioir industrial and technical capaboility. Yamaoto was for a time targeted for assasination. He was, however, appointed commander of the Imperial Navy. He was an inovative strategist and in particular propmoted the naval air wing. The League of Nations assigned the South Pacific Island Mandate (Nan-Yo) to Japan. Japan fought with the Allies in World War I. The Mandate covered a large part of Micronesia and were former German colonies. The Mandate included the Carolines, Marianas, Marshall Islands and Palau archepeligoes. An exception was made for Guam which continued under U. S. administration. The South Pacific Mandate was put under the control of the Japanese Navy. The Governors appointed were mostly admirals or vice-admirals. The Mandate capital was Koror, in the Palau islands. The most important island was Saipan in the Marianasa for both military and economic reasons. Sapan also became a major center for Japanese settlement. Another important island militarily was Truk (now Chuuk), in the Carolines. The Imperial Navy referred to it as the Southern Gibraltar. Japan withdrew from the League (1935). This invalidated the Mandate, but Japanese had by this time integrated the islands into the Japanese Empire. Large numbers of Japanese immigrants had settled on the islands, especially Saipan. The Japanese Navy built airfields, fortifications, ports, and other military instaltions. Construction of military instalations was escalated after Japan withdrew from the League. These islands were to be major battlefields in the Pacific War. The major naval powers (America, Britain, France, Italy, and Japan) agreed to substantial limitations on their naval strength which at the time was measured in battleships. American Secretary of State, Charles Evans Hughes organized a conference to address the problem of spiraling naval expendidutres as a result of the naval arms race. Senator William E. Borah, Republican of Idaho, who had led the fight against American ratification of the Treaty of Versailles and participation in the League of Nations, strongly advocated efforts to limit the arms race. His efforts were not at first favored by the new Harding administration, but was eventually adopted as the Republican alternative to the Democrat's (Wilson's) policy of collective security through the League of Nations. The Confrence opened on Armistice Day 1921--a very meaningful date so close to World War I. The American delegation was led by Secretary of State Charles Evans Hughes. Hughes shocked the other delegates by proposing a major reduction in naval fleets and not just limitations on new construction. This was far beyond what the other countries had anticipated. Some have called this one of the most dramatic moments in American diplomatic history. The American proposals entailed scrapping almost 2 million tons of warships as well as alengthy �holiday� on new building. The consequences of the Washington Treaties went far beyond this. The world naval powers convened in London to discuss continued naval arms limitations. The London Conference was strongly promoted by British Prime Minister Ramsey MacDonald who desired to continue if not increase the limitations established by the Washington Naval Treaties (1921). The Conference was held as the Wall Street Crash (1929)was spiraling into a serious world-wide economic criis and the major powers desired to cut back on gobernment sopending, especially military power. A Treaty was signed (April 1930). The signatories agreed to build no replacements of capial ships before 1937. American, Britain, and Japan agreed to avoid a arms race in destoyers and submarines. They also for the first time placed limits ob cruisers. America and Britain were allocated a cruisr tonnage about one and half that of Japan. The partipants agreed to another naval arms conferemnce in 1935. The inferior status of Japan has caused considerable resentment after the Washington Naval Conference (1921-22). After the London Conference it set in motion political changes of serious consequences. The Americans and the British attempted to convene another naval arms conference (1935). The major naval powers met in London for another round of naval talks to renew the existing limitations decided on at the Washington Naval Conference (1921-22) and London Naval Conference (1930). These limits were due to expire (1935-37). The militarsts in Japan were now in virtual control of the Government. The Japanese demanded parity with America and Britain. When this was not granted, the Japanese withdrew from the planned conference. This meant the exisiting limitations would expire. All three nations initiated battleship rebuilding programs with expiration of the treaty in 1936. Japan initiated the largest building program, a massive program to build 150 ships. The Japanese laid down two super battleships, Yamoto and Musashi, but the actual dimensions of these massive ships were kept secret. They were 69,100 tons, twice the size of treaty limitations. Germany built Bismarck and Tirpitz at 52,600 tons. The falure of the Conference created enough concen in Congress to approve an American naval building program, although a smaller program than initiated by the Japanese, only 100 vessels. Even so the new ships would only bring the Navy up Treaty limits. Two aircraft carriers were laid down in 1936 and 1937, each within Treaty limits. (These were USS Wasp (CV-7) and the larger USS Hornet (CV-8). No one knew at the time just how importnt these carriers would be. Both would reach the fleet in 1941 in time to participate in the critical Pavific battles of 1942. The Rooevelt Administration justified the appropriations in part as they would create jobs. The Isolationists and peace lobby opposed the appropriations with the slogan "Schools, not battleships". New battleships were authorized, but actual keels were not laid until after the war began in Europe. Only the USS North Carolina (BB-55) reached the fleet before Pearl Harbor. Yamamoto and other naval officers appreciated the industrial potential of America. Japanese Army commanders had no such apprecition. In addition, the Japanese Army in 1939 fought an undeclared war with bthe Soviets along the Manchurian border. The Japanese Army had suffered substantial losses and were not anxious to persure another campaign against the Soviets. American support for China caused Army officers to advocate a war with America. Many in the Army had convinced themselves that fighting spirit could over come American industrial superiority. As the British were engafed in Europe and France and the Netherlands occupied, their colonies with key natural resources needed by resource-poor Japan seem easy prey. The only other creditable force in the Pacific was the American Pacific Fleet at Pearl Harbor. The Army in 1941 dominated the Japanese Government. The U-boat successes in 1940-41 had a significant impact on the Pacific naval campaign. The Battle of the Atlantic was perhaps the most important single struggle in World War II. Defeat here would have forced Britain to capitulate and Aamerica could have not poarticipated in the European conflict. The power ballance in the Pacific was made even more lopsided by the fact that most of the Royal Navy had to be deployed in the Atlantic to keep the sealanes open. Until Bismarck was destroyed, the British battleships and carriers had to be concentrated in the North atlabtic. Even though America was legally neutral, by 1941 the Roosevlt Administration was not only supporting Britain, but engaging in an undeclared naval war in the North Atlantic. The Japanese Imperial Navy (Nihon Kaigun) by 1941 was the dominant naval force in the Pacific. The Japanese had a large well-trained naby with excellent ships. There was no peace-time neglect as the British and U.S. navies experienced. The American Navy was aware that the Japanese had a modern effective navy, but did not fully understand the campabilities of the Imperal Navy or the danger posed by the sizeable carrier fleet. The Japanese Imperial fleet was superbly trained with outstanding night fighting capabilities. Not only was the fleet well trained and included modern vessels, many of the Japanese vessels and naval aircraft were supperior to American and British vessels in many aspects. Japan led the world in operational aircraft carriers and carrier aircraft. (The British in 1941 were still using Swordfish biplanes on their carriers and American planes, especially the fighters were slower and less manuerable. The Japanese Mitsubishi Type 00 fighter, the Zero, was both faster and more maneuverable than either the U.S. Navy carrier fighter, the Grumman F4F Wildcat. The full extent of the threat was in part obscured by American rascial sterotypes and wide-spread belief that the Imperial Navy was not an effective force. In reterospect, the only suprising question about the attack on Pear Harbor and Japanese offensive in the South Pscifiuc, is not how they succedded, but how America managed to stop the Japanese after only 6 months of victories. The two glaring weakeneses were the lack of radar and the ineffective fore supression systems. Not well understood is that Japan had a very substantial submarine fleet. Looking back as a historian, it is almost incomprehensible that Japan decided to wage war against the United States. War with Britain and the Netherlands is more understandable. Britain in 1941 looked like if not a defeated nation, at least a severely weakened one. The Netherlands was occupied by Axis ally NAZI Germany. America is a very different matter. The United States was not at war. It had not been weakened by the War. And Japan had no commitment that the Germany would join them if they attacked America. War with America seems like an extrodinarily wreckless decession for a country already mired down in a war with China and that had experienced a sharp defeat in a short war with the Soviets. Why would Japan have decided on war with America, a country with a larger population and a much larger industrial and scientific base. The road to war began early in the history of modern Japan. Wars with China (1895), Russia (1904-05), and Germany (1914-18) proved both short and profitable, enabling Japan to build a small empire. The risring influence of the military brought to power men of limited outlook who saw military action as a legitimate use of sate power. They were backward looking men who saw the European empires of the 19th century as to what Japan should seek to establish. And they were men who were strongly influenced by the historic image of the Samurai and Bushido which convinced them that Japanese racial superiority and martial spirit could prevail over the material supperority of America. Despite the power of American induistry, tgey saw Americans as a weak, decadent people who would not fight. Most of the Japanese militarists who made this judgement on which the very life of Japan would hang, knew no Americans and had little or no experience with Ameica. It was the Japanese carrier attack on Pearl Harbor that brought America into the War. While Pearl Harbor was a stunning tactical victory, it was a strategic blunder by the Japanese of incaluable proportions. It was a stunningly successful military success, brilliantly executed by the Japanese. Eight battle ships, the heart of the American Pacific fleet were sunk. But the three carriers were not at Pearl. Despite the success of the attack, it was perhaps the greatest strtegic blunder in the history of warfare. The Japanese attack on the Pacific fleet at Pearl Harbor changed everything. A diverse and quareling nation, strongly pacifistic was instantly changed into a single united people with a burning desire to wage war. The issolationism that President Roosevelt had struggled against for over 7 years instantly disappeared. Even Lindburg asked for a commision to fight for the United States. Both the American and Japanese Navies had concluded before the War that if war was to come, it would be decided by a major fleet action at sea in which the deciding factor would be the big-gun battleships. The destruction of the American battleships at Pearl Harbor forced Admiral Nimitz when he took command of the Pacific Fleet to develop a new strategy. The only strategy open to Nimitz was to put the air craft carriers which had survived the Pear Harbor attack. The problem here was that the Japanese had a vastly sup[erior carrier force. Not only did the Japanese have many more carriers, but their pilots were more expeienced and better trained and the planes had superior performance characteristics. And the battle groups formed around the carriers after Pearl Harbor included cruisers to defend the carriers, vessels that were no match to the battleships of the Imperial Navy. The question was not wether the Navy could stop Japanese invasions of the Phillipines, Malaysia, and the Dutch West Indies, but wether the Pacific Fleet could even prevent its own destruction and defend the Hawaian Islands. With the American fleet impobilized at Pear Harbor, the Japanese were able to sweep through the Pacific and Southeast Asia. Guam was quickly taken. Resistance at Eake sland suprised the Japanese, but after the initial assault was repulsed, a second assault took the island. MacArthur's defense of the Philippines was compromised when most of his planes were destroyed on the fround at Clarke Field. General MacArthur commanded the most important American military force west of Pearl. His handlong of the defense of the Philippines wasdisapponting at best, bordering on incompetence. He failed to strike back at the Japanese in the hours after the attack on Pearl Harbor by bombing Jpanese bases in Formosa. He also allowed much of the available aircraft to be destroyed on the ground. [Schom] The horror of the Batan Death March created an impage of the Japanese military in the American mind that fueled a hatred for the Japanese. [Schom] Hong Kong quickly fell. The Japanese also seized the oil-rich Dutch East Indies (modern Indonesia). Allied naval forces fought a series of engagements to stop the Japanese, but could not match the powerful Japanese naval forces. Animitz and Halsey tried to distract the Japanese with hit an run carrier raids. The Japanese moved south from IndoChina, seizing Malayia and then the bastion at Singapore. The Repulse and Prince of Wales are lost in the defense of Singapore. Then they moved west through Thailand and defeating the British in Burma. Within a few months the Japanese had carved out the huge empire with enormous resources that they had long coveted. The Japnese then targeted New Guinea in preparation for a move south to Australia. All that remained to stop them were four American carriers. The Japanese strategic concept was to seize a huge empire and then fortify it so that it would be enormously costly for the Americans to retake. The resources from the empire which the Japanese called the Greater East Asia Co-Prosperity Sphere were to be used to support the Japanese military. The Japanese with little knowledge of America were convinved that America would never make the sacrifices needed to retake the Japanese conquests. This strategic concept was fataly flawed. First, the attack on Pearl Harbor turned a biterly divided America into a unified, mortal enemy. Second, the Japanese strategy had no provision for attacking the industrial base of the United States, an industrial base far exceeding the industrial capacity of Japan. Third, the Japanese were unprepared for the American submarine campaign, a campaign which by 1943 was beginning to deny Japanese industry the resources from their newly won empire. THe Japanese found their army bogged down in unwinnable campaigns in China and Burma and morooned on isolated Pacific islands that they could no longer supply. The radar developed by British and American scientists proved a huge advantage to the american Navy in the Pacific campaign. The Imperial Navy had no radar. This was one of many examples of the Axis failure to share military technoology. The Germans developed higly effective radar systemns, both for air an naval warfare. The Bismarck sunk by the British in 1941 (months before Pearl Harbor) for example had an effective radar system. If the Germans had supplied this technology to the Japanese in 1939-40, the Pacific naval campaign would have beren much more difficult than it was. The Japanese would, for example, fared much better at Midway. After the War had turned agasinst the Germany, a technology excjange was arranged with Japan. I am not sure radar was involved, but by then it was too late for the Imperial Navy which after the Battle of Letye Gulf was no longer an effective force. Allied radar and many other technical advances were the result of close cooperation between American and British scientists anf joint development projects that began even before America entered the War. There was no comparable Axis technical cooperation or even coordination of military campaigns. The Kriegsmarine had very effective radar on its surface ships like Bismarck yet advanced German technology like radar, jet engines, and other equipment was not provided to the Japanese until very late in the War, too late to be of any effectiveuse to the Japanese war effort. The Japanese after Pearl Harbor launched a series of invasions designed to seize territory that wouuld supply the natural resources that the miltarists so coveted. Here they were enormously successful. What they failed to do was to launch a knock out blow against the American Pacific fleet. Although the american carriers had not been at Pearl, the Japanese had a 3:1 superority in carriers. The Japanese superority was not just in numbers, but they had better planes and more experienced flight crews. Japan to win the War had to use its naval superority in 1942 to destroy the Pacific Fleet. Given America's industrial might, the naval superority would be rapidly closed by 1943. Thus Japan should have forced the Pacific fleet to battle in which it could deploy its massive carrier force against the Pacific Fleet's much smaller force. Instead of striking at Pearl, the Japanese instead deployed carriers to the Indian Ocean and in operations preparing for an invasion of Australia. Each of these were marginal undertakings in a war against America and the key American force, the Pacific Fleet. Japan failed to do this. In fact, the Pacific Fleet managed to deploy its carriers so effectively that the two major fleet engagements during 1942 (Coral Sea and Midyay) were fought on relaively equal terms. It is a mistake of clossal proprtions when you have a 3:1 superority in carriers to flight fleet engagements on equal terms. The primary maxim of carrier warfare is the force which first spots and launces on enemy carriers is likeky to win the engagement. For this reason you never want to fight an engagement on equal terms. With a larger carier force, you can launch on the enemy carriers even if he finds some of your carriers force. The Japanese ignored these basic maxims and as a result, the Imperial Navy's carrier superority was squandered and America had the time it needed to build a fleet of unimaginable size and power. The American Pacific fleet carriers by a fortuitous accident of history were not at Pearl wheen the Japanese struck (December 7), And it was the Pacific fleet carriers that were the primary Japanese targets. It is a massive Japanese failure. Admiral Halsey brought Enterprise into Pearl to the scene of a bloodied and still smoldering fleet (December 8). The ship was huredly resupplied and left Pearl to seek out the Japanese fleet (December 9). The Japanese carrier task force had long since departed. This was actually another fortuitous circumstance. If Enterprise or the other American carriers had encountered the six carriers of the Japanese task force, the resulting battle probably would have been disaterous. The Japanese at this point in the War were not only better trained, but flying superior aircraft types. And the Americans still were not able to coordinate multi-carrier attacks. (This defincy was still apparent at Midway.) Enterpise aircraft did find and sink the Japanese subnmarine I-70. After determining that the Japanese had departed, Admiral Nimitz had to devise a war plan for the Pacific's fleet's only remaining substantial force--the carriers (Enterprise, Lexington, and Yorktown). Nimitz had to prepare opperational plans with the understanding that new carriers could not be built and reach the fleet until 1943. Loss of the carriers would nean the Imperial Navy would totally dominate the Pacific and even cut off Pearl. The result was a decession not to risk using the carriers in a massive naval battle with the Imperial Fleet, but rather to carry out a series of swift, small-scale hit-and-run raids in the South Pacific. These raids were designed to disrupt Japanese buildups in islands from which anoher attack on Pearl could be staged. In addition, it would give the carrirs and the air groups aboard them the opportunity to practive their skill and operational effectiveness before taking on the Imperial Fleet. The Japanese occupied with carving out their Southeast Asian empire gave the Americn carriers the time needed to become more effective fighting machines. The Enterprise battle group commanded by Admiral Halsey hit the Marshalls (Wo Chi and Kwajeleen) (February 21). Then Enterprise hits Wake Island and Marcus Island. These raids were not without risk, bringing carriers wthin the range of land-based aircraft. The Pacific fleet's carrirs, however, emerged from these operations much more effective and capable fighting ships. One of these least noted naval campaign was the Indian Ocean engaements during early 1942. Admiral Nagumo with the First Air Fleet entered the Indian Ocean with a force of five carriers and four fast battleships as well as cruisers and destroyers (March 26, 1942). The purpose appears to have been to support Army operations in Burma and escort a convoy to Rangoon and then strike the Btitish naval base in Ceylon (Sri Lanka) where the Royal Navy had been building a substantial naval force. Incredibly this was a larger carrier force than deployed two months later against Midway. The force succeeded in sinking the British light carrier HNS Hermes, two cruisers, and smaller ships. The Royal Navy was asonished with the power of the Japanese carrier force. At this stage of the War, the Japanese carrier aircraft were far superiir to the British carrier aircraft. After the engagement the Royal Navy retired from the eastern Indian Ocean. It is unclear what the value of this campsign was. At the time the only creditable threat to Japan was the badly mauled American Pacific fleet and its four priceless carriers. Any assessment of the military situation would suggest that Japan should have focused on bringing the Pavific fleet to battle to get at those carriers. It is unclear what the purpose of this powerful firce was. They could have seized Ceylon or even attacked British facilities in India. While Nagumo had considerable success against the Royal Navy force, the Royal Air Force from bases in Ceylon had down or damaged a substantial number of Japanese planes. Nagumo had dispersed the British threat, but the American Pacific fleet carriers were still a threat and the British had impaired the combat effectiveness of the First Air Fleet The news from the Pacific was an unrelenting series of disasters. America needed a victory. The only intact offensive force in the Pacific was Americais carriers. Army Air Corps pilot with B-25s trained for carrier take offs. The B-25 was a medium bomber never intended for carrier use. Carrier commander Afm. "Bull" Halsey led a taskforce made up of Hornet and Enterprise. It was a risky operation as it committed half of the Pacific fleet's carrier force to a very dangerous operation. The B-25s took off from Hornet. It was the first blow to the Japanese home islands. The raid was led by Lt. Col. Jimmy Doolittle. The physical damage was inconsequential, but the psychological impact was immense. Most of the Amrican aviators crash landed in China and were helped to reach saftey by Chinese Nationalist guerillas. The Japanese reprisals were savage. A estimated 0.5-0.7 million Chinese civilians were murdered. The Japanese Navy was so embarassed that they rushed forward Admiral Yamaoto's plans to complete the job left unfinished at Pearl. Yamamoto conceived a plan that would bring the desimated American Pacific Fleet to battle at Midway Island and use the Imperial Fleet's massive superority to destoy it. The American code breaking effort was designated "Mafic". Magic included breaking both the Japanese diplomatic code and the naval code. At the time of Peal Harbor the Americans were reading the diplomatic, but not the naval code. Although American code breakers did not prevent the attack on Pearl Harbor, as American code breakers cracked the naval code it played a key role in the 1942 naval engagements in which the Japanese held vastly superior naval forces. Commander Joe Rochefort played the critical role in the code breaking effort. [Schom] The results enable the U.S. Navy to turn the Japanese back in the Coral Sea and defeat the Japanese at Miday--ll accomplished with inferior naval resources. A side benefit of Magic was that knowledge of the diplomatic code enabled American code breakers to read transmissions from the Japanese Embassy in Berlin. This provided detailed reports about fighting on the Eastrern Front, information the Soviets were reluctant to supply. The first important Allied effort to stop the Japanese sweep through the Pacific occurred in the Coral Sea. The Japanese planned to seize Port Moreseby, completing their conquest of New Guinea and a smaller operation in the Solomons at Tulagi. Port Moresby would have provided a launching pad for an invasion of Australia itself. (At the time, most of the Australian Army was in North Africa fighting Rommel's Afrika Korps.) The Japanese landing force was escorted by the front-line carriers Shokaku and Zuikaku. The Japanese naval task force en route to seize Port Moresby was intercepted by an American carrier force, alerted by American code breakers. It was the first carrirer to carrier engagement in history. The Japanese launched an attack on the Americans, but found only a destroyer and oiler. In the meantime the Americans sank the Japanese light carrier Shoho (May 7). The next day the two carrier forces fought a major engagement. The Japanese succeeded in sinking Lexington and heavily damaging Yorktown (May 8). The Americans heavily damaged Shokaku and devestated the air crew of Zuikaku. The substantial Japanese pilot casualties was very signigicant. Despite the American losses, the Japanese invasion force turned back, the first major Japanese reversal of the War. The Japanese assessment of the battle was that not only was Lexington sunk, but that Yorktown was either sunk or so badly damaged that it could no longer be deployed. This affected planning for the Miday operation. The engagement appears to have convinced Japanese naval planners that the American carriers were no mach for the Japanese carriers. The Japanese failed to preceive that the American carriers effectively fought the battle or that the surprise appearance of the American carrier in the Coral Sea to oppose the invasion of Port Moresby resulted from American code breaking. It also meant that they had lost a carrier, and large numbers of planes and pilots. This effectively removed two front line carriers from the Japanese order of battle. This reduced the available carriers for the Midway operation. Combined with the British damage to the First Air Fleet in the Indian Ocean, Admiral Yamamoto had allowed their carrier forces to be significantlseriously weakened in operations of marginal importance. This was critical because if Japan was to win the War it had to be done in 1942 when they had overwealming superiority in the Pacific. If the War developed into a war of attrition, the far greater indusstrial resources of the United States would prevail. Admiral Yamamoto planned the Midway Opeation as a war-winning stroke. He asseembled the most powerful force in naval history, up to that time. It outclassed the Pacific Fleet in every ship type. What Yamamoto did not understand was that fighting the battle with only four of its first-line carriers and off Midway, the Americans would actually muster more aircraft than the Japanese. Thus Yamamoto was actually creating the circumstance in which the Americans could establish a locallized superiority in air power. Midway proved to be the turning point of the Pacific War. It is notable because it was the only major Allied victory in which the opposing forces were superior. Admiral Yamamoto was determined to bring the American Pacific fleet to battle before America's industrial might could redress the strategic ballance. Yamamoto reasoned that Midway was an assett of such importance that Nimitz would have to commit his few remaining assetts to defend it. The Japanese had many advantages. Unknon to them, however, surprise was not one of the advantages. The same American code breaking operation that had learned of the Port Moresby operation also warned Admiral Nimitz that the next target was Midway. Admiral Yamamoto was convinced that the remaining American carriers could be brought to battle and destroyed at Midway. The Japanese plans were based on achieving an element of surprise and on the fact that two American carriers had been destoyed in the Coral Sea, in fact the Yorktown, although heavily damaged had not been sunk. American code breakers had alerted the Americans to the Japanese plans. Admiral Nimitz positioned Enterprise and Hornet, along with the hastily patched up Yorktown northwest of Midway to ambush he Japanese. The American carrier victory at Midway dealt a crippling blow to the Imperial Navy. The Americans sank four first-line Japnese carriers, killing most of the well-trained crews. The weakness of the Japanese in fire saftey and fire supression was notable. While the Imperial Navy still held an advantage, it was no longer an overwealming one. Meanwhile American shipyards were turning out the new Essex class carriers that would engage the weakened Imperial Navy in 1943. The stunning American carrier victory at Midway, significantly reduced the strike capability of the Imperial Navy. The Japanese had in a remarkably short period of time seized most of the territory that they had set out to obtain with the resiurces they coveted. And they had done so with relatively limited losses. Holding that territory and supplying the military units posted to hold and exploit the expansive new empire, however, proved to be a much more difficult undertaking. This required both naval vessels and a large maritime (maru) fleet. The Japanese maru fleet was adequate for peace time, but not for war and servicing their expabded new empire. The disaster at Midway significantly changed the military balance. It meant that further offensives would be contested with comparable forces and that existing positions were now open to attack. The Imperial Navy, however, did not aprise the Army of the dimensions of the Midway disaster. And Navy planners still believed they had the capability of destoying the American Pacific Fleet. Most of the Japanese Army was deployed in China. Units could be redeployed to the Pacific, but transporting and supplying them would further stress the capabiity of the Japanese maru fleet. One author describes the Japanese dilema, "The Imperial Navy had 132 Pacific island bastions assigned to the Yokosuka South Seas District, 27 stations assigned to the Kure District, and 23 to the Sasebo District." And all this did not include the forces bogged down in China and Manchuria as well as recntly committed to Malaya, Burma, and Indo-Chinathe, Philipppines, Borneo, and Dutch East Indies. All these forces had to be supplied by the maritime (maru) fleet. Japan did not have sufficet tonage to do this. [Jersey, p.9.] And this was even before the American submarine campaign had begun to effectvely target Japanese shipping. Major World War II battles in Europe were massive undertakings. The Pacific War was different. Major battle were fought by a handfull of divisions. Given the siuze of the islands and the limitations of supplying these island bastion, the Japanese were limited in the forces they could deploy. Thus major engagements which would determine the fate of the Empire, such as the battle for the Mrinas, were fought out by a relatively small force. After Midway the focus of the War shifted south. America had reduced the Japanese naval advantage, but did not yet have the naval assetts needed to chalenge the Imperal Fleet in a major fleet action in the Central Pacific. The Japanese with a badly-damaged Fleet Air Arm declined to renew challenge the America Pacific Fleet. Both sides instead began to regroup and rebuild their naval forces for a future show down in the Central Pacific. This was a serious mistake for the Japanese as time was on America's side. The tremendous industrial capacity of the United States could build naval vessels and aircraft at a far more rapid rate than Japan. Japan did renew its offense in the South Pacific which had been put on hold after the Coral Sea Battle. This was a natural development because the Japanese after taking the Dutch East Indies had seized almost all of New Guinea, except for Port Moresby and the island groups to the East, including New Britain, New Ireland, and the Solomons. Unlike the Central Pacific, these were large islands (especially New Guinnea) and lcated close together. Thus the fighting could be suported with air fields rather than carriers. Thus the fighting was largelky land operations intersperced with short range amphibious operations. The inintial phase of the campaign was Japanese assaults on Australia, bombing runs, a land offensive crossing the Owen Stanley Mountains to take Port Moresby, and building an air field on Guadacanal in the southern Solomons. This airfield could be used to support naval opertions to cut off Australia from American reinforcements and supplies. The center of the Japanese operations was the vast complex of military instalations the Japanese built at Rabaul on New Britains. The Japanese based their best pilots and planes there. The first American offensive of the Pacific War occurred when U,S, Marines seized the airfield the Japanese were building on Guadalcanal (August 1942). What followed was one of the most prolonged campaigns foughht by the Marines in the War and a series of pitched battles in in the Slot formed by the Solomon Islands. The subsequent Allied offensive was a two prong movement. The first prong was overseen by the U.S. Navy (Halsey) in the Solomons and other islands east of New Guinea. The Navy decided against a costly assault on Rabaul itself. Rather they established rings around Rabaul, cutting off the powerful base and making it impossible for the Japanese to resupply it. They subjected Rabaul to a whithering air assault. Allied troops on Los Negros in the Admiralty Islands played a major role in cutting off and neutralizing Rabaul (December 1943). The second prong was overseen by the U.S. Army (MacArthur) with Australian support. The Australians stopped the Japanese short of Port Moresby. American infantry began taking bases along the northern coast of New Guinea. MacArthur's goal from the beginning was to obtain bases from which he could return to Philippine Islands. Bases in New Zealand brought the southern Philippine Islands into range. For the first time in naval warfare, aircraft played an important role. A major aspect of the War was that the carrier replaced the battleships as the key capital ship. And the carrier was nothing more than a floating airfield capable of moving aircraft in range of enemy fleet formations and land targets. Naval aviation has two components, the planes and the pilot who flew them. Only three countries (America, Britain, and Japan) built and deployed carriers. The Germans had plans to do so, but military reverses prevented them from doing so. The carrier aircraft that fought in the Pacific are some of the most storied aircraft of World war II. The Japanese began the War with the most effective carrier aircraft, especially the elegant, but lightly armored A6M Mitsubishi Zero (1941). Not only did the Japanese lose carriers at the Coral Sea and Miday, but the core of their carrier pilots was desimated and further attrited in South Pacific campaign. The Japanese had a superb pilot training program, but it was highly selective and long. It was designed to produce small numbers of suberb pilots. As the short, quick war turned into an extended war of attrition, the Japanese did not modify their training program. They were thus unable to turn out competent carrier pilot to replace the extensive losses in 1942. The Japanese training program was also hamperd by oil shortages. America in contrast launched an extensive pilot training program which crewed the new cartriers flowing out of shipyards in incredible numbers. The American pilots were not as well trasined with the initial Japanese pilots, but they were competntly trained and soon gained battle expoerience. In addition a new generation of Americn planes reached the fleet which were superior to the lightly armoured Zero. The new Navy F4U Corsair and F6F Hellcat, combined with the Army Air Corps USAAF P-51 and P-47 fighters devestated Japanese planes. The Japanese fought the entire war with the same basic planes they began the War with, one of many examples of the limited industrial capability with whjivh the Japanese fought the War. The smll force of American carriers and epecially the victory at Midway provided the time for American industrial capacity to build a naval force with which Japan's limited industrial capacity could not cope. The Japanese built 10 additional by 1945, but the Americans built over 150 carriuers. American shipyards had a substantial capacity, but the American ability to produce ships of all descriptions not only astonished the Japanese and Germans, but it surprised American naval planners as well. The reason was the innovative ship building methods developed by Henry Kaiser and others to produce the Liberty Ships. While the German submarine campaign in the North Atlantic failed, the American submarine campaign in the Pacific proved spectacularly successful. Hampered by ineffective torpedoes in 1942, the American submarines by 1943 began to significantly affect the delivery of raw materials to Japan. The American submarines targeted the Japanese merchant marine (maru) fleet. While the big fleet carriers got the headlines. The American submarines sunk over 50 percent of all vessels destroyed during the War. The Japanese merchant marine was almost completely destroying, cutting the country's war industries off from supplies and bringing the country close to starvation by 1945. The American submarines did to Japan what the German u-boats tried to do to Britain. Surprisingly the Japanese submarine fleet had little impact on the Pacific campaign. Unlike the Americans, the Japanese began the War with the effective Type 93 Long-Lance Torpedo. The Japanese Navy never used their submarines to interdict American supply vessels. Rather they were used to target fighting ships with only limited success because of their tactical deployment. The Japanese used their submarines as scouts and to targer warships. As the American offensive moved toward the Home Islands, the Japanese used their submarines to supply bypassed island garisons, some of which were near starvation. They were also used to supply bypassed island bases where many garrisons were close to starvation. They also managed to get some secret German military technology to Japan late in the war (1944-45). While MacArthur and the Army drove fowad in New Guiena, the Navy achieved most of its objectives in the South Pacific/Solomons. Rabaul was not taken as had been planned with Operatioin Cartwheel, but it was neutralized. And the Imperial Navy's pilot forcehad been drained to support air operationd from Rabaul. The Navy next began a new offensive drive in the Central Pacific. Amercan industrial strength by now was producing ships and planes that enabled America to build a naval force capable of leap froging from island to island. The Japanese had gambled when they struck Pearl Harbor. They believed that Americans were morally weak and would not fight. Guadacanal and Tarawa showed that they were very wrong. And now they faced the most powerful industrial power in the World which they had roused to a nation demanding retribution. Worse still there German allies were falling back on all fronts. Japan would be left to fight both American and Britain alone. Unlike the Solomons campign, the Imperial Navy did not contest the American advance cross the Central Pacific. The Combined Fleet even withdrew from Truk, its primarit Central Pacific bastion. Only when the Americans assalted the Marianas did the Imperial Navy sally forth. The Marianas were the key to the War. The new American B-29 could reach the Japanese Home Islands from the Marianas and the Japanese knew it. The result would be the epic Battle of the Phillipnes Sea. The Imperial Navy had been preparing for the battle for 2 years. The result shocked the Empire to its core. The two American offensives in the Pacific came to a conclusion at the same time. The U.S. Army under Douglas MacArthur in the South Pacific had neutrilized Rabaul and defeated or bypassed Japanese forces in the Solomons and northeastern New Guinea. At the same time, the U.S. Navy under Admiral Chester Nimitz after driving through the Central Pacific (the Giberts and Marshalls) and finally seieed the Marianas after the great naval victory in the Philippines Sea. But this brought to the fore the still unanswered question of 'where next?' There were two targets on the table. MacArthur was adament about the answer--the Philippines. Since departing Corrigedor he had repeated his goal, 'I shall return.' His argument was largely political and moral--we owed it to the Filipino people as the Philippines at the time was American territory. Admiral Earnest King believed that Formosa (Taiwan) made more strategic sense, largely because it would more more effectively interdict the delivery of raw materials from the Southern Resourse Zone to the Home Islands. A difference of such magnitude between such senior American commanders could only be resolved by President Roosevelt. The President summoned his commanders at Pearl Harbor, Hawaii to settle the issue of the direction of the advance on Japan (July 26-27). MacArthur made his and the Army's case. Nimitz made the case for the Navy. The choice would be the Philippines leading to the greatest naval battle in world history--the Battle of Leyte Gulf. The Allies at the Second Quebec Conference/Octagon (September 12-16) approved plans fo a British Pacific Fleet (BPF) to join the Americans in the Pacific War as the War in Europe was winding down. The Conference was attendened by President Roosevelt, Prime-Minister Churchill, and the Joint Chiefs. Admiral King was not entusiastic, but the President supported Churchill. The major contribution was to be the formation of the BCP. At the time there was hope that the War would be over by Christmas. As this proved overly optimistic, the hidt of forces east was delayed. The British were also thinking about a bomber force to operate from Luzon and Okinawa. It was to be called Tiger Force. Gen. Arnold unlike Admoral King supported the idea. There were also plans to redeploy Commonwealth forces for the invasion of Japan. The Americans were not all that interestec at Quebec, but after the firece fighting for Iwo Jima and Okinawa, they became more interested. [Weinberg, p. 844.] The bomber Tiger Force and a British infantry invasion force was never formed. The BPF was and actually saw action off the Home Islands. The BPF was made up of Commonwealth naval vessels. The BPF oficially ame into being (November 22) from the ships of the former British Eastern/East Indies Fleet. It was initially based at Trincomalee, Ceylon (Sri Lanka). Once constituted, the base was moved to Sydney, Australia and a forward base on Manus Island, part of the Admiralty Chain north of New Guinea. Manus was taken back from the Japanese (February-March 1944). It was a substantial force, one of the largest fleets ever assembled by the Royal Navy. It consisted of 4 battleships, 6 fleet aircraft carriers, 15 smaller aircraft carriers, 11 cruisers, and numerous smaller warships, submarines, and support vessels. While the U.S. Navy was not initially enthusiastic, had the atomic bombs not ended the War, the BPF would have been needed for the invasion of Japan. The Battle of Leyte Gulf was the greatest naval battle in history. It was the last engagement where opposing battleships faced each other. It was a desperate action by Japan which essentially pitted its battleships against American carriers. Japan had begun the war with a stratgy based in carrier warfare, but by 1944 the Imperial Navy was firced to rely on its battleships against an overwealming american carrier force. The Battle for Leyte Gulf occurred October 23-26, 1944. The Japanese drew up a highly complex strategy to throw virtually the entire remaining Imperial fleet in a desperate attempt to oppse the American invasion of the Philippine Islands. The battle evolved in four separate actions thayt both side found difficult to coordinate in the furious battles that ensyed. Once the American landings at Leyte Gulf began the Japanese ordered three separate forces to oppose the Americans. The Japanese Cental Force or main force sailed through the Philippines to reach the American landings at Leyte. It was spotted and intercepted by American carrier from Task Force 38. It was mauled in the Sibuyan Sea. The giant Nustachi was sunk. It turned back, but was not destroyed. The Japanese Southern Force tried tried to reach Leyte through the Surigao Strait. They were intercepted by American battleships, some of which had been raised from Pearl Harbor. This proved to be the greatest surface action since Jutland in World War I and the last important action between battleships in history. American destroyers and battleships using radar in a night action virtually wiped out the Southern Force. Next the Northern Force which was a carrier force meant to decoy Halsey's Third Fleet was spotted. Halsey immediately raced to engage this force. It was descimated in an engagement off Cape Engano. Meanwhile the Japanese Central Force had turned around and was approaching Leyte Gulf. Unprotected by the Third Fleet, the Japanese would have wreaked hhavic on the troop ships and supply vessels. The only force between Leyte Gulf was a small group of American escort carriers and destroyers. The esort carriets were supporting the landings and protecting the invasion fleet from submine attack. They did not have armour piercing bombs needed for attacks on battleships. Somehow this force managed to turn around the Japanse Central Force in a an action off Samar, although at great cost. After the battle the Imperial Fleet no longer existed as a creditable naval force. [Thomas] From the onset of the Pacific War, the South China Sea becme essentially a Japanese lake. Japanese merchant vessels brought the resources of the Southeast Asia (Southern Resource Area) through the South China Sea without fear of American interdiction. This only bgan to change in late 1943 when the American sunmarine campaign began to resolve the problems encountered early n the War, but Allied surface ships still could not operate in the South China Sea through 1944. This finally changed in 1945, in part because of the deterioration of Japanese naval and air power after the Battle of Leyte Gulf and in invasion of the Philippines (October 1944). Task Force (38) was one of the most powerfull naval forces ever assembled. It was tasked with supporting the Luzon landings at Lingayen. As part of this operation it was tasked with engaging Japanese air fields and naval targets to the east capable of supporting Japanese forces in the Philippines. TF 38 also swept the South China Sea. The retaking of Clark Field and other air bases on the Philippines meant that American carrier groups were now able to enter the South China Sea. TF 38 began with raids on airfields on Formosa (January 3) and for a week hammered other targets. TF 38 entered the South China Sea through the Bashi Channel between Formosa and Luzon (January 10). TF 38 was composed of both Essex Fleet cariers and lighter carriers and support vessels. TF 38 hit Saigon and Tourane Bay, Indochina (January 12). TF 38 destroyed 44 enemy ships including 15 naval vessels as well as many aircraft and military facilities. TF 38 struck Formosa and Canton in China (January 15). Further strikes hit Canton again and Hong Kong (January 16). TF 38 exited the South China Sea through the Balintang Channel, another channel, between Taiwan and Luzon (January 20). They headed north, striking targets on Formosa (January 21) and Okinawa (January 22). TF 38 finally reached Ulithi Lagoon for resupply and refitting (January 26). The Navy South China Sea operation along with the seizue of the Phillipnes and the American submarine campaign meant that Japan was cut off from their Southern Resource Area. Japan launched the War to obtain the resources of Southeast Asia and now those resources were no longer available to the war industries on the Home Islands. Okinawa was a major land campaign conducted by the Army and Marines. The ininital landings were unopposed, but the Japanese made their stand in the south and the fighting was brutal. The Japanese exacted a terrible toll from carefully prepared defenses. The American landings were largely unopposed by the Imperial Navy which had been desimated at the Battle of Leyte Gulf. The Imperial Navy was no longer capable of a fleet action to oppose the American Navy. The great naval battles of the War were now over. The primary exception was the suiside mission of the Yamato, the remaining Japanese super-battleship. The Yamato had turned around at Samar. This time there would be no turning back. The Yamato was intercepted by American carrier aircraft and sunk after repated hits in a horific bombardment. The Navy's role. however, was critical. Carrier aircraft was needed to provide air support until the land needed for air bases could be seized from the Japanese and airfilds established. Also the Navy was needed for logistical support, landing the assault troops and then supplying them. This meant the Navy had to bring its ships into range of Japanese lnd-based aircraft. The Japanese mounted a major Kamakazee campaign against the American ships off Okinawa. The Japanese aircraft were for the most part the same aircraft that the Japanese had begun the war with. Most of the well-trained pilots were gone. Japan would conduct its Kamakazee campaign with obsolete air craft and virtually untrained pilots, butlaunched in mass, some could get through. The Kamakazee was the lone effective weapon that could strike at the American invasion fleet. The success on Okinawa was fully understood by the Japanese and was to be the basis for the defense of the Home Islands. They assembled a secret air force of 5,000 planes to be used in a massive Kamikazzee attack on the llied invasion fleet. The Yamato sally was oart of the Jaoanese response to th American invasion of Okinawa, it would prove to be the final important naval action of the Pacific War. Yamato and her sister ship Musashi were built in secret. Buildings were built to obstruct views of the shipyards. Like Bismarck and Tirpitz they exceeded naval arms limitations. The Japanese ships were evern larger than the German golithas. They were built to win what both the American and Japanese expected to settle the Pacific War--a climatic bug0gun battleship engagement. The Japanese realized that they could not outbuild the Americans in number. So like the Germans they opted for quality. The two Japanese battleshios were the largest ever built. They had phenomenal 18 in guns. The modern American battle ships onlky had 16 inch guns. Thiscgave them a longer range, although not much longer. Range was important because it maean that the Japanese could engage American battkle ships before the Americans could return fire. Of course this was all moot because by the time these ships were on station, the American the carriers could engage Japamese ships at huge differnce with bombs and torpedoes. The Japanese held back both 'Yamato' and 'Musashi' for the climatic battle. Unfortunately for the Imperial Fleet, the Japamese began losing battleship in the Solomoms (late-1942), long before the climatic battle that they were planning. They finally committed their fleet in the Battle of Leyte Gulf (October 1944). They lost Musachi amd Yamato was damaged. Worse in the Imperial Navy's mind, Adm. Kondo turned the fleet around in thr face of the enemy as he was nearing Leyte Gulf and the American transports. This at a time that the Arny was engagong in mass suiside banazi charges and Kamikazee flights. The Yamato sally was aimed at regaining the Navy's honor. The Americams begam to land on Okinawa (April 1). They were surprised that there was no resistance at the beach. Gen. Ushijima realised that he had no manser to American naval artillery. He would fight inland from carefilly prepared defenses anf Kamikaze suiside air attacks. The Imperial Navy to remitself from its Ammiral Kondo's withdrawl from Leyte Gulf decided on a spectaculr suiside attack of its own. The Yamato had been repaired and enough fuel crped toether for a one way sally. Japanese cadets at the Etajima Academy (Hiroshima) were given the honor of helping to crew Yamato. American Marines and Army soldiers cut the island in two (April 4). The first important Japanese response came in the form of a 350-gun Kamikaze suiside raid (April 6-7). Part of theis effort was the Yamato The plan was after the aerial Kamikaze attacks for Yamato to finish off what was left of the American fleet and thn having used up its fuel, beach itself and provide artillery support to Ushijima's 32nd Army. The Yamato's fate would be the dramatic end of big-gun battleships. American pre-War naval planing envisioned a climatic naval action with big gun battleships to win a war in the Pacific. Immediately after Pearl Harbor, the United States attempted to launch a commerce war against the Japanese. Japan like Britain was an industrial island natiin which had to import both food and raw materials. Oil was particularlyb important because there is almost no domestic production. One might ask why a nation so vulnerable woukd launch a naval war against an industrial giant, but the Japanese obsession with winning the war in China and seizing the Southern Resource Zone (SRZ) was so intense tht the Japanese military decided to gamble. A commerce war against Japanese merchant shipping (marus) was an obvious naval strategy. And Japan had begun the War with an inadequate merchant fleet. The demand of supplying its nely seized, far-fling empire put a grear srain on the fleet. Admiral Nimitz was an old submariune commander and the submarines were one of few naval ssetts available. Faulty torpedoes, however, rendered early sunmarine actions ineffective. This problem was not completely resolved until 1943 and it is at this time tht the Americans submriners began to seriously impact the maru fleet. They were assisted by Ultra intercepts. The destruction of the Imperial Navy in the Battles of the Philippines Sea and Leyte Gulf left the Japanese powerless to keep sea lane open except for coastal shipping. Japan had the needed resources in the SRZ to continue the War, but no way of getting those resources to the Home Islands. It is at this time that the U.S. Navy surface fllet and aircraft in firt the Philippins and then Okinawa joined the submarines tightened the blockade of the Home Islands. The next step was the Strategic boming campaign which hammered Japanese ports. Carrier aircraft and land based planes from Iwo Jima and Okinawa commenced low-level attacks which hit small craft including fishing boars and barges. The final step was dropping aerial bombs which left all the major ports inperable even if ships slipped through the blockade. This completed the blockade and left the Japanese despertely short of food and facing a poor 1945 harvests. Civilians were put to work collecting acorns. Food rations were cut to subsistence levels. The few war plants that somehow survived the bombing were brought to a hault because of the shortage of raw materials. Some of the last reserves of petroleum were used to create the fuel needed for a last desperate Kakikazze attack aimed at the Allied invasion fleet. There was not even sufficent fuel to properly train new pilots. America ended the warby dropping two atomic bombs were dropped (August 1945) and Japan finally surrendered (September 1945).The American Manhattan Program was initiated by President Roosevelt when work done by German physicists led to concern that the NAZIs might build an atomic bomb. Jewish and oher refugees fleeing the NAZIs made a major contribution to the success of the Manhattan Program. The first bomb was successflly tested at Alamagordo, New Mexico on July ??, 1945. The Allies met in a Berlin suburb after the NAZI surrender to make dcisions about the occupation of Germany and defeating Japan. The Allied powers 2 weeks after the bomb was tested demanded on July 27, 1945 that Japan surrender unconditionally, or warned of "prompt or utter destruction". This became known as the Potsdam Declaration. The Japanese military was prepared to fight on rather than surender. The Japanese Government responded to the Potsdam Declaration with "utter contemp". The Japanese military continued feverish pland to repel the Ameican invasion of the Home Islands. Many Whermacht generals at the end of the War were anxious to surrnder to the Amreicans. One German General commanding forces west of Berlin after the War said, "We wondered why they didn't come." This was not the attitude of the Japanese military. I know of know memoir written by an important Japanese military officer expresing similar sntiments. Truman was not anxious to use the atomic bomb. He was anxious to end the War and limit Ameican casulties. For Truman the Japanese response to the Potsdam Declaration made up his mind. There have been many books and aticles published in both Japan and America about the atomic bomb. Japanese scholars have reserched the decission making process that led to the dropping of the atomic bombs. Almost always the focus is on Truman and Ameican military leasers. Rarely do Japanese authors address the role of Japanese political and military leaders. The United States dropped two atomic bombs on Hiroshima and Nagasaki on August 6 and 9, and the Soviet Union entered the war against Japan on August 8. Freidel, Frank. Franklin D. Roosevelt: Rendezuous with Destiny (Little Brown: Boston, 1990), 710p. Jersey, Stanley Coleman. Hell's Island: The Untold Story of Guadalcanal (Texas A&m: College Station, 2008), 514p. Schom, Alan. The Eagle and the Rising Sun: The Japanese-American War 1941-1943 (Norton, 2003). Thomas, Evan. Sea of Thunder: Four Commanders and the Last Great Naval Campaign, 1941-1945 (Simon & Schuster: New York, 2006), 414p. Weinberg, Gerhard l. A World at Arms: A Global History of Wlorld War II (Cambrige Universit Press: New York, 2005), 1178p. 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Altruism or selflessness is the principle or practice of concern for the welfare of others. It is a traditional virtue in many cultures and a core aspect of various religious traditions and secular worldviews, though the concept of "others" toward whom concern should be directed can vary among cultures and religions. Altruism or selflessness is the opposite of selfishness. Altruism can be distinguished from feelings of loyalty. Pure altruism consists of sacrificing something for someone other than the self (e.g. sacrificing time, energy or possessions) with no expectation of any compensation or benefits, either direct, or indirect (e.g., receiving recognition for the act of giving). Much debate exists as to whether "true" altruism is possible. The theory of psychological egoism suggests that no act of sharing, helping or sacrificing can be described as truly altruistic, as the actor may receive an intrinsic reward in the form of personal gratification. The validity of this argument depends on whether intrinsic rewards qualify as "benefits." The term altruism may also refer to an ethical doctrine that claims that individuals are morally obliged to benefit others. Used in this sense, it is usually contrasted with egoism, which is defined as acting to the benefit of one's self. - 1 The notion of altruism - 2 Individual variations - 3 Scientific viewpoints - 4 Religious viewpoints - 5 Philosophy - 6 See also - 7 Notes - 8 References - 9 External links The notion of altruism The concept has a long history in philosophical and ethical thought. The term was originally coined in the 19th century by the founding sociologist and philosopher of science, Auguste Comte, and has become a major topic for psychologists (especially evolutionary psychology researchers), evolutionary biologists, and ethologists. Whilst ideas about altruism from one field can have an impact on the other fields, the different methods and focuses of these fields always lead to different perspectives on altruism. In simple terms, altruism is caring about the welfare of other people and acting to help them. A certain individual may behave altruistically in one case and egoistically in another situation. However, some individuals tend to behave more altruistically, while others tend to behave more egoistically. Altruism may be considered a general attitude to the point where altruism has been considered as a trait. A 1986 study estimated that altruism was half-inherited.[notes 1] Another study, by the Massachusetts Institute of Technology, was published in 2009. Also based on the twin design, the new study estimated that genetic differences accounted for approximately 20% of individual variations.[notes 2] Marcel Mauss's book The Gift contains a passage: "Note on alms." This note describes the evolution of the notion of alms (and by extension of altruism) from the notion of sacrifice. In it, he writes: Alms are the fruits of a moral notion of the gift and of fortune on the one hand, and of a notion of sacrifice, on the other. Generosity is an obligation, because Nemesis avenges the poor and the gods for the superabundance of happiness and wealth of certain people who should rid themselves of it. This is the ancient morality of the gift, which has become a principle of justice. The gods and the spirits accept that the share of wealth and happiness that has been offered to them and had beenhitherto destroyed in useless sacrifices should serve the poor and children. - Compare Altruism (ethics) – perception of altruism as self-sacrifice. - Compare explanation of alms in various scriptures. In the science of ethology (the study of animal behaviour), and more generally in the study of social evolution, altruism refers to behaviour by an individual that increases the fitness of another individual while decreasing the fitness of the actor. In evolutionary psychology this may be applied to a wide range of human behaviors such as charity, emergency aid, help to coalition partners, tipping, courtship gifts, production of public goods, and environmentalism. Theories of apparently altruistic behavior were accelerated by the need to produce theories compatible with evolutionary origins. Two related strands of research on altruism have emerged from traditional evolutionary analyses and from evolutionary game theory a mathematical model and analysis of behavioural strategies. Some of the proposed mechanisms are: - Kin selection. That animals and humans are more altruistic towards close kin than to distant kin and non-kin has been confirmed in numerous studies across many different cultures. Even subtle cues indicating kinship may unconsciously increase altruistic behavior. One kinship cue is facial resemblance. One study found that slightly altering photographs so that they more closely resembled the faces of study participants increased the trust the participants expressed regarding depicted persons. Another cue is having the same family name, especially if rare, and this has been found to increase helping behavior. Another study found more cooperative behavior the greater the number the perceived kin in a group. Using kinship terms in political speeches increased audience agreement with the speaker in one study. This effect was especially strong for firstborns, who are typically close to their families. - Vested interests. People are likely to suffer if their friends, allies, and similar social ingroups suffer or even disappear. Helping such group members may therefore eventually benefit the altruist. Making ingroup membership more noticeable increases cooperativeness. Extreme self-sacrifice towards the ingroup may be adaptive if a hostile outgroup threatens to kill the entire ingroup. - Reciprocal altruism. See also Reciprocity (evolution). - Direct reciprocity. Research shows that it can be beneficial to help others if there is a chance that they can and will reciprocate the help. The effective tit for tat strategy is one game theoretic example. Many people seem to be following a similar strategy by cooperating if and only if others cooperate in return. - One consequence is that people are more cooperative if it is more likely that individuals will interact again in the future. People tend to be less cooperative if they perceive that the frequency of helpers in the population is lower. They tend to help less if they see non-cooperativeness by others and this effect tend to be stronger than the opposite effect of seeing cooperative behaviors. Simply changing the cooperative framing of a proposal may increase cooperativeness such as calling it a "Community Game" instead of a "Wall Street Game." - A tendency towards reciprocity implies that people will feel obligated to respond if someone helps them. This has been used by charities that give small gifts to potential donors hoping thereby to induce reciprocity. Another method is to announce publicly that someone has given a large donation. The tendency to reciprocate can even generalize so people become more helpful toward others in general after being helped. On the other hand, people will avoid or even retaliate against those perceived not to be cooperating. People sometimes mistakenly fail to help when they intended to, or their helping may not be noticed, which may cause unintended conflicts. As such, it may be an optimal strategy to be slightly forgiving of and have a slightly generous interpretation of non-cooperation. - People are more likely to cooperate on a task if they can communicate with one another first. This may be due to better assessments of cooperativeness or due to exchange of promises. They are more cooperative if they can gradually build trust, instead of being asked to give extensive help immediately. Direct reciprocity and cooperation in a group can be increased by changing the focus and incentives from intra-group competition to larger scale competitions such as between groups or against the general population. Thus, giving grades and promotions based only on an individual's performance relative to a small local group, as is common, may reduce cooperative behaviors in the group. - Indirect reciprocity. The avoidance of poor reciprocators and cheaters causes a person's reputation to become very important. A person with a good reputation for reciprocity have a higher chance of receiving help even from persons they have had no direct interactions with previously. - Strong reciprocity. A form of reciprocity where some individuals seem to spend more resources on cooperating and punishing than would be most beneficial as predicted by several established theories of altruism. A number of theories have been proposed as explanations as well as criticisms regarding its existence. - Pseudo-reciprocity. An organism behaves altruistically and the recipient does not reciprocate but has an increased chance of acting in a way that is selfish but also as a byproduct benefits the altruist. - Costly signaling and the handicap principle. Since altruism takes away resources from the altruist it can be an "honest signal" of resource availability and the abilities needed to gather resources. This may signal to others that the altruist is a valuable potential partner. It may also be a signal of interactive and cooperative intentions since those not interacting further in the future gain nothing from the costly signaling. It is unclear if costly signaling can indicate a long-term cooperative personality but people have increased trust for those who help. Costly signaling is pointless if everyone has the same traits, resources, and cooperative intentions but become a potentially more important signal if the population increasingly varies on these characteristics. - Hunters widely sharing the meat has been seen as a costly signal of ability and research has found that good hunters have higher reproductive success and more adulterous relations even if they themselves receive no more of the hunted meat than anyone else. Similarly, holding large feasts and giving large donations has been seen as ways of demonstrating one's resources. Heroic risk-taking has also been interpreted as a costly signal of ability. - Both indirect reciprocity and costly signaling depend on the value of reputation and tend to make similar predictions. One is that people will be more helping when they know that their helping behavior will be communicated to people they will interact with later, is publicly announced, is discussed, or is simply being observed by someone else. This have been documented in many studies. The effect is sensitive to subtle cues such as people being more helpful when there were stylized eyespots instead of a logo on a computer screen. Weak reputational cues such as eyespots may become unimportant if there are stronger cues present and may lose their effect with continued exposure unless reinforced with real reputational effects. Public displays such as public weeping for dead celebrities and participation in demonstrations may be influenced by a desire to be seen as altruistic. People who know that they are publicly monitored sometimes even wastefully donate money they know are not needed by recipient which may be because of reputational concerns. - Women have been found to find altruistic men to be attractive partners. When looking for a long-term partner more conventional altruism may be preferred which may indicate that he is also willing to share resources with her and her children while when looking for a short-term partner heroic risk-taking, which may be costly signal showing good genes, may be more preferable. Men also perform more altruistic acts in the early stages of a romantic relationship or simply when in the presence of an attractive woman. While both sexes state that kindness is the most preferable trait in a partner there is some evidence that men place less value on this than women and that women may not be more altruistic in presence of an attractive man. Men may even avoid altruistic women in short-term relationships which may be because they expect less success. - People may compete over getting the benefits of a high reputation which may cause competitive altruism. On other hand, in some experiments a proportion of people do not seem to care about reputation and they do not help more even if this is conspicuous. This may possibly be due to reasons such as psychopathy or that they are so attractive that they need not be seen to be altruistic. The reputational benefits of altruism occur in the future as compared to the immediate costs of altruism in the present. While humans and other organisms generally place less value on future costs/benefits as compared to those in the present, some have shorter time horizons than others and these people tend to be less cooperative. - Explicit extrinsic rewards and punishments have been found to sometimes actually have the opposite effect on behaviors compared to intrinsic rewards. This may be because such extrinsic, top-down incentives may replace (partially or in whole) intrinsic and reputational incentives, motivating the person to focus on obtaining the extrinsic rewards, which overall may make the behaviors less desirable. Another effect is that people would like altruism to be due to a personality characteristic rather than due to overt reputational concerns and simply pointing out that there are reputational benefits of an action may actually reduce them. This may possibly be used as derogatory tactic against altruists, especially by those who are non-cooperators. A counterargument is that doing good due to reputational concerns is better than doing no good at all. - Group selection. It has controversially been argued by some evolutionary scientists such as E. O. Wilson that natural selection can act at the level of non-kin groups to produce adaptations that benefit a non-kin group even if these adaptions are detrimental at the individual level. Thus, while altruistic persons may under some circumstances be outcompeted by less altruistic persons at the individual level, according to group selection theory the opposite may occur at the group level where groups consisting of the more altruistic persons may outcompete groups consisting of the less altruistic persons. Such altruism may only extend to ingroup members while there may instead prejudice and antagonism against outgroup members (See also in-group favoritism). Group selection theory has been criticized by many other evolutionary scientists. Such explanations do not imply that humans are always consciously calculating how to increase their inclusive fitness when they are doing altruistic acts. Instead, evolution has shaped psychological mechanisms, such as emotions, that promote altruistic behaviors. Every single instance of altruistic behavior need not always increase inclusive fitness; altruistic behaviors would have been selected for if such behaviors on average increased inclusive fitness in the ancestral environment. This need not imply that on average 50% or more of altruistic acts were beneficial for the altruist in the ancestral environment; if the benefits from helping the right person were very high it would be beneficial to err on the side of caution and usually be altruistic even if in most cases there were no benefits. The benefits for the altruist may be increased and the costs reduced by being more altruistic towards certain groups. Research has found that people are more altruistic to kin than to no-kin, to friends than to strangers, to those attractive than to those unattractive, to non-competitors than to competitors, and to members ingroups than to members of outgroup. The study of altruism was the initial impetus behind George R. Price's development of the Price equation, which is a mathematical equation used to study genetic evolution. An interesting example of altruism is found in the cellular slime moulds, such as Dictyostelium mucoroides. These protists live as individual amoebae until starved, at which point they aggregate and form a multicellular fruiting body in which some cells sacrifice themselves to promote the survival of other cells in the fruiting body. Selective investment theory proposes that close social bonds, and associated emotional, cognitive, and neurohormonal mechanisms, evolved in order to facilitate long-term, high-cost altruism between those closely depending on one another for survival and reproductive success. Such cooperative behaviors have sometimes been seen as arguments for left-wing politics such by the Russian zoologist and anarchist Peter Kropotkin in his 1902 book Mutual Aid: A Factor of Evolution and Peter Singer in his book A Darwinian Left. Most recently, Jeremy Griffith has proposed a biological theory for the development of truly altruistic instincts that accommodates the biological imperative to reproduce, as evidenced by a moral conscience visible in humans today. Jorge Moll and Jordan Grafman, neuroscientists at the National Institutes of Health and LABS-D'Or Hospital Network (J.M.) provided the first evidence for the neural bases of altruistic giving in normal healthy volunteers, using functional magnetic resonance imaging. In their research, published in the Proceedings of the National Academy of Sciences USA in October 2006, they showed that both pure monetary rewards and charitable donations activated the mesolimbic reward pathway, a primitive part of the brain that usually lights up in response to food and sex. However, when volunteers generously placed the interests of others before their own by making charitable donations, another brain circuit was selectively activated: the subgenual cortex/septal region. These structures are intimately related to social attachment and bonding in other species. Altruism, the experiment suggested, was not a superior moral faculty that suppresses basic selfish urges but rather was basic to the brain, hard-wired and pleasurable. Another experiment funded by the National Institutes of Health and conducted in 2007 at the Duke University in Durham, North Carolina suggests a different view, "that altruistic behavior may originate from how people view the world rather than how they act in it." In the study published in the February 2007 print issue of Nature Neuroscience, researchers have found a part of the brain that behaves differently for altruistic and selfish people. The researchers invited 45 volunteers to play a computer game and also to watch the computer play the game. In some rounds, the game resulted in the volunteers winning money for themselves, and in others it resulted in money being donated to a charity of the volunteer's choice. During these activities, the researchers took functional magnetic resonance imaging (fMRI) scans of the participants' brains and were "surprised by the results". Although they "were expecting to see activity in the brain's reward centers," based on the idea that "people perform altruistic acts because they feel good about it," what they found was that "another part of the brain was also involved, and it was quite sensitive to the difference between doing something for personal gain and doing it for someone else's gain." That part of the brain is called the posterior superior temporal cortex (pSTC). In the next stage, the scientists asked the participants some questions about type and frequency of their altruistic or helping behaviours. They then analysed the responses to generate an estimate of a person's tendency to act altruistically and compared each person's level of altruism against their fMRI brain scan. The results showed that pSTC activity rose in proportion to a person's self-reported level of altruism. According to the researchers, the results suggest that altruistic behavior may originate from how people view the world rather than how they act in it. "We believe that the ability to perceive other people's actions as meaningful is critical for altruism," said lead study investigator Dharol Tankersley. The International Encyclopedia of the Social Sciences defines psychological altruism as "a motivational state with the goal of increasing another’s welfare." Psychological altruism is contrasted with psychological egoism, which refers to the motivation to increase one’s own welfare. There has been some debate on whether or not humans are truly capable of psychological altruism. Some definitions specify a self-sacrificial nature to altruism and a lack of external rewards for altruistic behaviors. However, because altruism ultimately benefits the self in many cases, the selflessness of altruistic acts is brought to question. The social exchange theory postulates that altruism only exists when benefits outweigh costs. Daniel Batson is a psychologist who examined this question and argues against the social exchange theory. He identified four major motives for altruism: altruism to ultimately benefit the self (egoism), to ultimately benefit the other person (altruism), to benefit a group (collectivism), or to uphold a moral principle (principlism). Altruism that ultimately serves selfish gains is thus differentiated from selfless altruism, but the general conclusion has been that empathy-induced altruism can be genuinely selfless. The empathy-altruism hypothesis basically states that psychological altruism does exist and is evoked by the empathic desire to help someone who is suffering. Feelings of empathic concern are contrasted with feelings of personal distress, which compel people to reduce their own unpleasant emotions. People with empathic concern help others in distress even when exposure to the situation could be easily avoided, whereas those lacking in empathic concern avoid helping unless it is difficult or impossible to avoid exposure to another's suffering. Helping behavior is seen in humans at about two years old, when a toddler is capable of understanding subtle emotional cues. In psychological research on altruism, studies often observe altruism as demonstrated through prosocial behaviors such as helping, comforting, sharing, cooperation, philanthropy, and community service. Research has found that people are most likely to help if they recognize that a person is in need and feel personal responsibility for reducing the person's distress. Research also suggests that the number of bystanders witnessing distress or suffering affects the likelihood of helping (the Bystander effect). Greater numbers of bystanders decrease individual feelings of responsibility. However, a witness with a high level of empathic concern is likely to assume personal responsibility entirely regardless of the number of bystanders. A feeling of personal responsibility or - moral norm - has also strongly been associated with other pro-social behaviors such as charitable giving. Many studies have observed the effects of volunteerism (as a form of altruism) on happiness and health and have consistently found a strong connection between volunteerism and current and future health and well-being. In a study of older adults, those who volunteered were significantly higher on life satisfaction and will to live, and significantly lower in depression, anxiety, and somatization. Volunteerism and helping behavior have not only been shown to improve mental health, but physical health and longevity as well. One study examined the physical health of mothers who volunteered over a 30-year period and found that 52% of those who did not belong to a volunteer organization experienced a major illness while only 36% of those who did volunteer experienced one. A study on adults ages 55+ found that during the four-year study period, people who volunteered for two or more organizations had a 63% lower likelihood of dying. After controlling for prior health status, it was determined that volunteerism accounted for a 44% reduction in mortality. Merely being aware of kindness in oneself and others is also associated with greater well-being. A study that asked participants to count each act of kindness they performed for one week significantly enhanced their subjective happiness. It is important to note that, while research supports the idea that altruistic acts bring about happiness, it has also been found to work in the opposite direction—that happier people are also kinder. The relationship between altruistic behavior and happiness is bidirectional. Studies have found that generosity increases linearly from sad to happy affective states. Studies have also been careful to note that feeling over-taxed by the needs of others has conversely negative effects on health and happiness. For example, one study on volunteerism found that feeling overwhelmed by others' demands had an even stronger negative effect on mental health than helping had a positive one (although positive effects were still significant). Additionally, while generous acts make people feel good about themselves, it is also important for people to appreciate the kindness they receive from others. Studies suggest that gratitude goes hand-in-hand with kindness and is also very important for our well-being. A study on the relationship happiness to various character strengths showed that "a conscious focus on gratitude led to reductions in negative affect and increases in optimistic appraisals, positive affect, offering emotional support, sleep quality, and well-being.". Psychologists generally refer to this virtuous cycle of helping others, doing good and subsequently feeling good as "the helper's high". "Sociologists have long been concerned with how to build the good society" ("Altruism, Morality, and Social Solidarity". American Sociological Association.). The structure of our societies and how individuals come to exhibit charitable, philanthropic, and other pro-social, altruistic actions for the common good is a largely researched topic within the field. The American Sociology Association (ASA) acknowledges Public sociology saying, "The intrinsic scientific, policy, and public relevance of this field of investigation in helping to construct 'good societies' is unquestionable" ("Altruism, Morality, and Social Solidarity" ASA). This type of sociology seeks contributions that aid grassroots and theoretical understandings of what motivates altruism and how it is organized, and promotes an altruistic focus in order to benefit the world and people it studies. How altruism is framed, organized, carried out, and what motivates it at the group level is an area of focus that sociologists seek to investigate in order to contribute back to the groups it studies and "build the good society". Most, if not all, of the world's religions promote altruism as a very important moral value. Buddhism, Christianity, Hinduism, Islam, Jainism, Judaism and Sikhism, etc., place particular emphasis on altruistic morality. Altruism figures prominently in Buddhism. Love and compassion are components of all forms of Buddhism, and are focused on all beings equally: love is the wish that all beings be happy, and compassion is the wish that all beings be free from suffering. "Many illnesses can be cured by the one medicine of love and compassion. These qualities are the ultimate source of human happiness, and the need for them lies at the very core of our being" (Dalai Lama). Since "all beings" includes the individual, love and compassion in Buddhism are outside the opposition between self and other. It is even said that the distinction between self and other is part of the root cause of our suffering. In practical terms, however, since most of us are spontaneously self-centered, Buddhism encourages us to focus love and compassion on others, and thus can be characterized as "altruistic." Many would agree with the Dalai Lama that Buddhism as a religion is kindness toward others. Still, the notion of altruism is modified in such a world-view, since the belief is that such a practice promotes our own happiness: "The more we care for the happiness of others, the greater our own sense of well-being becomes" (Dalai Lama). In the context of larger ethical discussions on moral action and judgment, Buddhism is characterized by the belief that negative (unhappy) consequences of our actions derive not from punishment or correction based on moral judgment, but from the law of karma, which functions like a natural law of cause and effect. A simple illustration of such cause and effect is the case of experiencing the effects of what I cause: if I cause suffering, then as a natural consequence I will experience suffering; if I cause happiness, then as a natural consequence I will experience happiness. In Buddhism, karma (Pāli kamma) is strictly distinguished from vipāka, meaning "fruit" or "result". Karma is categorized within the group or groups of cause (Pāli hetu) in the chain of cause and effect, where it comprises the elements of "volitional activities" (Pali sankhara) and "action" (Pali bhava). Any action is understood to create "seeds" in the mind that sprout into the appropriate results (Pāli vipaka) when they meet the right conditions. Most types of karmas, with good or bad results, will keep one in the wheel of samsāra; others will liberate one to nirvāna. Buddhism relates karma directly to motives behind an action. Motivation usually makes the difference between "good" and "bad", but motivation also includes the aspect of ignorance; so a well-intended action from an ignorant mind can easily be "bad" in that it creates unpleasant results for the "actor." In Buddhism, karma is not the only cause of all that happens. As taught in the early texts, the commentarial tradition classified causal mechanisms governing the universe in five categories, known as Niyama Dhammas: - Kamma Niyama — Consequences of one's actions - Utu Niyama — Seasonal changes and climate - Biija Niyama — Laws of heredity - Citta Niyama — Will of mind - Dhamma Niyama — Nature's tendency to produce a perfect type The fundamental principles of Jainism revolve around the concept of altruism, not only for humans but for all sentient beings. Jainism preaches the view of Ahimsa – to live and let live, thereby not harming sentient beings, i.e. uncompromising reverence for all life. It also considers all living things to be equal. The first Thirthankar, Rishabh introduced the concept of altruism for all living beings, from extending knowledge and experience to others to donation, giving oneself up for others, non-violence and compassion for all living things. Jainism prescribes a path of non-violence to progress the soul to this ultimate goal. Jains believe that to attain enlightenment and ultimately liberation, one must practice the following ethical principles (major vows) in thought, speech and action. The degree to which these principles are practiced is different for householders and monks. They are: - Non-violence (Ahimsa); - Truthfulness (Satya); - Non-stealing (Asteya); - Celibacy (Brahmacharya); - Non-possession or non-materialism (Aparigraha); A major characteristic of Jain belief is the emphasis on the consequences of not only physical but also mental behaviors. One's unconquered mind with anger, pride (ego), deceit, greed and uncontrolled sense organs are the powerful enemies of humans. Anger spoils good relations, pride destroys humility, deceit destroys peace and greed destroys everything. Jainism recommends conquering anger by forgiveness, pride (ego) by humility, deceit by straight-forwardness and greed by contentment. The principle of non-violence seeks to minimize karmas which limit the capabilities of the soul. Jainism views every soul as worthy of respect because it has the potential to become Siddha (Param-atma – "highest soul"). Because all living beings possess a soul, great care and awareness is essential in one's actions. Jainism emphasizes the equality of all life, advocating harmlessness towards all, whether the creatures are great or small. This policy extends even to microscopic organisms. Jainism acknowledges that every person has different capabilities and capacities to practice and therefore accepts different levels of compliance for ascetics and householders. The "great vows" (mahavrata) are prescribed for monks and "limited vows" (anuvrata) are prescribed for householders. In other words, the house-holders are encouraged to practice the five cardinal principles of non-violence, truthfulness, non-stealing, celibacy and non-possessiveness with their current practical limitations while the monks have to observe them very strictly. With consistent practice, it will be possible to overcome the limitations gradually, accelerating the spiritual progress. Altruism is central to the teachings of Jesus found in the Gospel, especially in the Sermon on the Mount and the Sermon on the Plain. From biblical to medieval Christian traditions, tensions between self-affirmation and other-regard were sometimes discussed under the heading of "disinterested love", as in the Pauline phrase "love seeks not its own interests." In his book Indoctrination and Self-deception, Roderick Hindery tries to shed light on these tensions by contrasting them with impostors of authentic self-affirmation and altruism, by analysis of other-regard within creative individuation of the self, and by contrasting love for the few with love for the many. Love confirms others in their freedom, shuns propaganda and masks, assures others of its presence, and is ultimately confirmed not by mere declarations from others, but by each person's experience and practice from within. As in practical arts, the presence and meaning of love becomes validated and grasped not by words and reflections alone, but in the making of the connection. St Thomas Aquinas interprets 'You should love your neighbour as yourself' as meaning that love for ourselves is the exemplar of love for others. Considering that "the love with which a man loves himself is the form and root of friendship" and quotes Aristotle that "the origin of friendly relations with others lies in our relations to ourselves," he concluded that though we are not bound to love others more than ourselves, we naturally seek the common good, the good of the whole, more than any private good, the good of a part. However, he thinks we should love God more than ourselves and our neighbours, and more than our bodily life—since the ultimate purpose of loving our neighbour is to share in eternal beatitude: a more desirable thing than bodily well being. In coining the word Altruism, as stated above, Comte was probably opposing this Thomistic doctrine, which is present in some theological schools within Catholicism. Many biblical authors draw a strong connection between love of others and love of God. 1 John 4 states that for one to love God one must love his fellowman, and that hatred of one's fellowman is the same as hatred of God. Thomas Jay Oord has argued in several books that altruism is but one possible form of love. An altruistic action is not always a loving action. Oord defines altruism as acting for the other's good, and he agrees with feminists who note that sometimes love requires acting for one's own good when the other's demands undermine overall well-being. German philosopher Max Scheler distinguishes two ways in which the strong can help the weak. One way is a sincere expression of Christian love, "motivated by a powerful feeling of security, strength, and inner salvation, of the invincible fullness of one’s own life and existence". Another way is merely "one of the many modern substitutes for love, ... nothing but the urge to turn away from oneself and to lose oneself in other people’s business." At its worst, Scheler says, "love for the small, the poor, the weak, and the oppressed is really disguised hatred, repressed envy, an impulse to detract, etc., directed against the opposite phenomena: wealth, strength, power, largesse." In Islam, the concept 'īthār' (altruism) is the notion of 'preferring others to oneself'. For Sufis, this means devotion to others through complete forgetfulness of one's own concerns, where concern for others is rooted to be a demand made by Allah on the human body, considered to be property of Allah alone. The importance lies in sacrifice for the sake of the greater good; Islam considers those practicing īthār as abiding by the highest degree of nobility. This is similar to the notion of chivalry, but unlike that European concept, in i'thar attention is focused on everything in existence. A constant concern for Allah (i.e. God) results in a careful attitude towards people, animals, and other things in this world. This concept was emphasized by Sufis of Islam like Rabia al-Adawiyya who paid attention to the difference between dedication to Allah (i.e. God) and dedication to people. Thirteenth-century Turkish Sufi poet Yunus Emre explained this philosophy as "Yaratılanı severiz, Yaratandan ötürü" or We love the creature, because of The Creator. For many Muslims, i'thar must be practiced as a religious obligation during specific Islamic holidays. However, i'thar is also still an Islamic ideal to which all Muslims should strive to adhere at all times. Judaism defines altruism as the desired goal of creation. The famous Rabbi Abraham Isaac Kook stated that love is the most important attribute in humanity. This is defined as bestowal, or giving, which is the intention of altruism. This can be altruism towards humanity that leads to altruism towards the creator or God. Kabbalah defines God as the force of giving in existence. Rabbi Moshe Chaim Luzzatto in particular focused on the 'purpose of creation' and how the will of God was to bring creation into perfection and adhesion with this upper force. Modern Kabbalah developed by Rabbi Yehuda Ashlag, in his writings about the future generation, focuses on how society could achieve an altruistic social framework. Ashlag proposed that such a framework is the purpose of creation, and everything that happens is to raise humanity to the level of altruism, love for one another. Ashlag focused on society and its relation to divinity. Altruism is essential to the Sikh religion. The central faith in Sikhism is that the greatest deed any one can do is to imbibe and live the godly qualities like love, affection, sacrifice, patience, harmony, truthfulness. The fifth Nanak, Guru Arjun Dev Sacrificed his life to uphold 22 Carats of pure truth, the greatest gift to humanity, the Guru Granth. The Ninth Nanak, Guru Tegh Bahadur Sacrificed his head to protect weak and defenseless people against atrocity. In the late 17th century, Guru Gobind Singh Ji (the tenth guru in Sikhism), was in war with the Moghul rulers to protect the people of different faiths, when a fellow Sikh, Bhai Kanhaiya, attended the troops of the enemy. He gave water to both friends and foes who were wounded on the battlefield. Some of the enemy began to fight again and some Sikh warriors were annoyed by Bhai Kanhaiya as he was helping their enemy. Sikh soldiers brought Bhai Kanhaiya before Guru Gobind Singh Ji, and complained of his action that they considered counterproductive to their struggle on the battlefield. "What were you doing, and why?" asked the Guru. "I was giving water to the wounded because I saw your face in all of them," replied Bhai Kanhaiya. The Guru responded, "Then you should also give them ointment to heal their wounds. You were practicing what you were coached in the house of the Guru." It was under the tutelage of the Guru that Bhai Kanhaiya subsequently founded a volunteer corps for altruism. This volunteer corps still to date is engaged in doing good to others and trains new volunteering recruits for doing the same. Advaita Vedanta differs from the view that karma is a law of cause and effect but instead additionally hold that karma is mediated by the will of a personal supreme god. This view of karma is in contradiction to Buddhism, Jainism and other Indian religions that do view karma as a law of cause and effect. Swami Sivananda, an Advaita scholar, reiterates the same views in his commentary synthesising Vedanta views on the Brahma Sutras, a Vedantic text. In his commentary on Chapter 3 of the Brahma Sutras, Sivananda notes that karma is insentient and short-lived, and ceases to exist as soon as a deed is executed. Hence, karma cannot bestow the fruits of actions at a future date according to one's merit. Furthermore, one cannot argue that karma generates apurva or punya, which gives fruit. Since apurva is non-sentient, it cannot act unless moved by an intelligent being such as a god. It cannot independently bestow reward or punishment. There exists a wide range of philosophical views on man's obligations or motivations to act altruistically. Proponents of ethical altruism maintain that individuals are morally obligated to act altruistically. The opposing view is ethical egoism, which maintains that moral agents should always act in their own self-interest. Both ethical altruism and ethical egoism contrast with utilitarianism, which is the view that every individual's well-being (including one's own) is of equal moral importance. A related concept in descriptive ethics is psychological egoism, the thesis that humans always act in their own self-interest and that true altruism is impossible. Rational egoism is the view that rationality consists in acting in one's self-interest (without specifying how this affects one's moral obligations). - Charity (practice) - Charitable organization - Comedy of the commons - Effective altruism - Family economics - Gene-centered view of evolution - Giving Pledge, pledge by Gates, Buffett and others to donate to charity at least half of their wealth - Inclusive fitness - Group selection - Kin selection - Mutual aid - The Power of Half; how a family came to decide to sell its home, so that it could donate half the proceeds to charity - Prisoner's dilemma - Prosocial behavior - Random act of kindness - Reciprocal altruism - Social psychology - Solidarity (sociology) - Tit for tat - The study covered 5 related traits, giving a broad heritability estimate of 56% for altruism, but then approximated all to 50%, and raised the maximum-likelihood heritability to 60% after corrections for unreliability. - The study generally refers to the trait studied as "giving". 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The wage-price spiral is a self-reinforcing process of wage increases and price increases that leads to inflation. This spiral occurs when workers demand higher wages in order to keep up with rising prices. In turn, firms raise prices in order to cover the higher cost of labor, leading to even higher wages and prices. As prices continue to rise, the spiral creates a feedback loop that drives inflation. The wage-price spiral is a key mechanism by which inflation can become self-sustaining. Once inflation starts to rise, the spiral can cause it to accelerate, leading to ever-higher prices. This can create a situation in which people expect prices to continue to rise, leading them to demand even higher wages. The wage-price spiral is a key factor in creating and sustaining inflation in an economy. What are the 3 main causes of inflation? The most common definition of inflation is an increase in the price level of a basket of goods and services over time. The main causes of inflation are: 1) Demand-pull inflation: This occurs when there is an increase in aggregate demand (AD) in the economy. This can be caused by a number of factors, including population growth, an increase in government spending, or an increase in investment spending. 2) Cost-push inflation: This occurs when there is an increase in the cost of production, such as an increase in the price of oil. This can cause a decrease in aggregate supply (AS), leading to higher prices. 3) Structural inflation: This occurs when there is an imbalance in the economy between aggregate demand and aggregate supply. This can be caused by a number of factors, including a change in tax rates, a change in government spending, or a change in the money supply. What does it mean when an economy faces inflation? Inflation occurs when there is an increase in the overall level of prices in an economy. This can be caused by a number of factors, including an increase in the money supply, a decrease in the production of goods and services, or a decrease in the purchasing power of currency. When inflation occurs, the value of currency decreases and the cost of living increases. This can lead to a decrease in purchasing power and a decrease in the standard of living. Why do wage increases cause inflation? Inflation is defined as an increase in the price level of goods and services. When wages increase, the cost of production also increases. This increase in cost is then passed on to the consumers in the form of higher prices. As the prices of goods and services increase, the purchasing power of the dollar decreases. This decrease in purchasing power is what we call inflation. What is inflationary spiral in economics? Inflationary spiral is a situation where inflation increases continuously, leading to higher prices for goods and services. This situation is often caused by an increase in money supply, which leads to higher prices for goods and services. What is the most common cause of inflation? The most common cause of inflation is an increase in the money supply. This can happen when the government prints more money or when the central bank increases the money supply through quantitative easing.
Number System is a way to represent the numbers in the computer architecture. There are four different types of number system such as binary number system (base 2), octal number system (base 8), decimal number system(base 10), and hexadecimal number system (base 16). In this article, let us discuss what is a binary number system, conversion from one system to other systems, table, positions, binary operations such as addition, subtraction, multiplication, and division, uses and solved examples in detail. Table of Contents: - Decimal to Binary Number Conversion - Binary Arithmetic Operations What is a Binary Number System? Binary Number System: According to digital electronics and mathematics, a binary number is defined as a number that is expressed in the binary system or base 2 numeral system. It describes numeric values by two separate symbols; basically 1 (one) and 0 (zero). The base-2 system is the positional notation with 2 as a radix. The binary system is applied internally by almost all latest computers and computer-based devices because of its direct implementation in electronic circuits using logic gates. Every digit is referred to as a bit. Binary Number System Table Some of the binary notations of decimal numbers are mentioned in the below list. \|Decimal number\|\|Binary value\| What is Bit in Binary Number? \|A single binary digit is called a “Bit”. The binary number below has 6 bits. Binary numbers chart with decimal values is organized in the below column. To understand binary math, you should know how the number system works. Let’s start with the decimal system since it is easier to deal with. Decimal to Binary Number Conversion For example, the number to be operated is 1235. 1235 = 1 × 1000 + 2 × 100 + 3 × 10 + 5 × 1 \|1000\|\|= 103 = 10 × 10 × 10\| \|100\|\|= 102 = 10 × 10\| \|10\|\|= 101 = 10\| \|1\|\|= 100 (any value to the exponent zero is one)\| The above table can be described as, 1235 = 1 × 1000 + 2 × 100 + 3 × 10 + 5 × 1 = 1 × 103 + 2 × 102 + 3 × 101 + 5 × 100 The decimal number system operates in base 10 wherein the digits 0-9 represent numbers. In binary system operates in base 2 and the digits 0-1 represent numbers and the base is known as radix. Put differently, the above table can also be shown in the following manner. We place the digits in columns 100, 101 and so on in base 10. When there is a need to put a value higher than 9 in the form of 10(n+1) for instance, to add 10 to column 100, you need to add 1 to the column 101. We place the digits in columns 20, 21 and so on in base 2. To place a value that is higher than 1 in 2n, you need to add 2(n+1). For instance, to add 3 to column 20, you need to add 1 to column 21. Position in Binary Number System In the Binary system, we have ones, twos, fours etc… For example 1011.110 It is shown like this: 1 × 8 + 0 × 4 + 1 × 2 + 1 + 1 × ½ + 1 × ¼ + 0 × 1⁄8 = 11.75 in Decimal To show the values greater than or less than one, the numbers can be placed to the left or right of the point. For 10.1, 10 is a whole number on the left side of the decimal and as we move more left, the number place gets bigger (Twice). The first digit on the right is always Halves ½ and as we move more right, the number gets smaller (half as big). In the example given above: - “10” shows ‘2’ in decimal. - “.1” shows ‘half’. - So, “10.1” in binary is 2.5 in decimal. Binary Arithmetic Operations Adding two binary numbers will give us a binary number itself. This is the simplest method. Addition of two single-digit binary number is given in the table below. \|1\|\|1\|\|0; Carry →1\| Let us take an example of two binary numbers and add them. Subtracting two binary numbers will give us a binary number itself. This is also an easy method. Subtraction of two single-digit binary number is given in the table below. \|0\|\|1\|\|1; Borrow 1\| Let us take an example of two binary numbers and subtract them. Example: Subtract 11012 and 10102. The multiplication process is the same for the binary numbers as it is for numerals. Let us understand it with example. Example: Multiply 11012 and 10102. The binary division is similar to the decimal number division method. We will learn with an example here. Example: Divide 10102 by 102 Uses of Binary Number System Binary numbers are commonly used in computer applications. All the coding and languages in computers such as C, C++, Java. etc. uses binary digits 0 and 1 to write a program or encode any digital data. Basically, the computer understands only coded language. Therefore these 2-digit number system is used to represent a set of data or information in discrete bits of information. Binary Number System Examples Let us practice some of the problems for better understanding: Question 1: What is binary number 1.1 in decimal? Step 1: 1 on the left-hand side is on the one’s position, so it’s 1. Step 2: The one on the right-hand side is in halves, so it’s 1 × ½ Step 3: so, 1.1 = 1.5 in decimal. Question 2: Write 10.112 in Decimal? 10.11 = 1 x (2)1 + 0 (2)0 + 1 (½)1 + 1(½)2 = 2 + 0 + ½ + ½ So, 10.11 is 2.75 in Decimal. Keep visiting BYJU’S to explore and learn more such math-related topics in a fun and engaging way.
What is SQL? SQL stands for “Structured Query Language”. It is a domain specific language which means it can be used to play with Relational Database only. As the name suggests SQL is used in handling structured data. Structured data is the data that has some specific structure/ format. It can be stored in tables like format or flat files likes CSV and TSV. It is a language used to query tabular data. It is an ANSI standard language used for manipulating, storing and accessing data in a database. This is a standardized query language for processing data stored in RDBMS (Relational Database Management System). Before understanding SQL we need to understand RDBMS first. RDBMS or Relational Database Management System is simply a database that stores structured data and there is a relation between the data. It stores data in tabular format. It has column and rows which contains related data entries. Columns are a vertical entity of a table. It contains the attribute of records and rows are the horizontal entity which contains records/data. The intersection of rows and column contains the information of a record with respect to that attribute. Given below is an example of a table. It is a language used to query over tabular data. Unlike other languages, SQL is a declarative language, one just needs to specify the result that they want to see and submit the query to RDBMS. RDBMS executes the code at the backend and gives the desired output. Whereas in a procedural language we have to tell a computer each and every step to perform in order to get the output. So if you want to select data from the above table you just need to write below query and execute it. SELECT * FROM CustomerDetail; A confusion with SQL is the syntax of SQL query. The elements are not executed in the order they are used in the query. Consider selecting data from above CustomerDetail table. SELECT ID, Name, Age FROM CustomerDetail WHERE Age > 20 Order By ID DESC; The above query will select all the records which have an age greater than 20 and display the result by order of ID. The sequence of execution of elements are as follows : - FROM: In the query FROM clause is executed first. It selects the tables and joins tables to get the base data. - WHERE: This clause filters the base data. So that there are fewer records in further processing. - GROUP BY: Group By clause combines rows into groups to perform aggregation. - HAVING: This clause is used to filter the aggregated data on the basis of the calculated column. - SELECT: This clause returns the selected records in the format requested by the user. - ORDER BY: This clause sorts the final data. So the lexical order and logical order of clauses in a SQL query differ but one needs to take care of these things when the performance comes into the picture. For smaller data retrieval user has to just mention the output he/she expects. How does SQL make working so easy The most important feature of SQL which makes it easy to work with is that it hides the complexity of processing. Since it is a declarative language, the programmer just needs to specify the format of output as per the requirement and the server will take care of all the complexity of retrieval and aggregation. So the code to retrieve data from a table will be smaller if written in SQL as compared to code written in any other language. It deals only with database objects. This is an advantage as well as a limitation of SQL. Because of this, it can be used to handle only structured data. with limited objects and structured data, working is easy in SQL. Even after dealing only with structural data it has more importance than any other programming language and it is easy to learn. This is based on basic relational algebra and tuple calculus. It takes just a few days to learn the basics of SQL. One can also learn this from an online tutorial. But becoming an expert and getting performance related expertise is an altogether different thing in SQL. It will take some time and hands-on experience. It also supports all the mathematical and string functions to modify the data according to need. It has all the features provided in any other programming language. This makes it an easier language to work with. Every programming language requires to interact with back end database and this has extensibility that it can be integrated into any language. Thus making it easy to work with any other programming language. 4.5 (1,730 ratings) Top SQL Companies Almost every IT company uses a database to store its data and manage it. But big companies that has a large set of data to deal with are the best to explore data. Given below are some of the top IT companies using SQL: - Tech Mahindra Various subset of SQL SQL queries can be categories into 4 main Category: 1. DDL (Data Definition Language) As the name suggests these types of queries are used to define the structure of data. Like the structure of a table, schema and modify it. Example – - CREATE: This command is used to create tables, database, schema etc. - DROP: This command is used to drop tables and other database objects. - ALTER: This command is used to alter the definition of database objects. - TRUNCATE: This command is used to remove tables, procedures, views, and other database objects. - ADD COLUMN: This command is used to add any column to table schema. - DROP COLUMN: This command is used to drop a column from any table structure. 2. DML (Data Manipulation Language) This type of queries is used to manipulate data in the database. Example – - SELECT INTO: This command is used to select data from one table and insert into another table. - INSERT: This command is used to insert data/records into a table. - DELETE: This command is used to delete records from the table. - UPDATE: This command is used to update the value of any record in the database. 3. DCL (Data Control Language) This category of SQL queries deals with the access rights and permission control of the database. Example – - GRANT: This command is used to grant access rights on database objects. - REVOKE: This command is used to withdraw permission from database objects. 4. TCL (Transaction Control Language) The transaction is a set of commands that perform a specific task on objects in a single unit of execution. So TCL commands deals with transactions in a database. Example – - COMMIT: This command is used to commits a transaction. Once committed it cannot be rolled back. This means the previous image of the database before running this transaction cannot be retrieved. - ROLLBACK: Rollback is used to revert the steps in transactions if an error occurs. - SAVEPOINT: This command sets a savepoint in the transaction to which steps can be rolled back. - SET TRANSACTION: This command is used to set characteristics of the transaction. What can you do with SQL It is mainly used in SQL SERVER MANAGEMENT STUDIO, a tool to manage database and data. It was launched by Microsoft for configuring, managing and administrating all the components of the database. Given below are the main operation one can do with SQL: 1. Create Database It can be used to create Database and it’s other objects. One can create a table to store data, stored procedure, functions to process data and views to view data. The user can also play around with joining data from different tables and get meaningful output. 2. Access Database A user can also manage the access rights on the database and its objects using SQL. One can check which user has executed which query and also privilege user has. An administrator can grant and revoke access from a user. 3. Manage Database Managing data is not an easy task. Especially when it’s important to business and has a huge size. So efficient storage and retrieval of data are important. SQL lets you do that without any hassle. 4. Manipulating Database These commands help you manipulate your data. Insert data into tables, delete records, update records all can be done easily using SQL commands. A user can also join different tables and have a view on collective data. 5. Website Use This can also be used with the integration of another programming language. Every programming language has an extension to embed SQL in its code. Working with SQL As we know this is a querying language, and it deals with the data stored in the back end. Hence the interface is not so interesting. One won’t get UI to play with colors and designs. There are just tables with columns and rows. But if data really interests you then SQL is the language you must learn. Working with SQL you get to play with data, join tables and performance tuning. You can write some procedure and transactions to perform analysis task and also schedule a job using SQL. Advantages of SQL Below are some of the advantages of SQL 1. Requires no coding This is a declarative language, one just needs to mention the output he/she wants. It has straight forward commands to perform actions like a select, update, delete etc. One does not need to write complex code to retrieve data from a database or manipulate the data. 2. Well defined standard It is an ANSI standard language. It has been established as a standard language for querying RDBMS. 3. Interactive Language It is used to communicate with a database and its objects. We can get the output of complex queries within seconds. 4. Manipulating Database It’s easy to update records in SQL and maintain the integrity of data. The relationship can also be implemented between two tables. It can be integrated with other languages to connect with database. The SQL query can be embedded in any other programming language used for application development. Required SQL Skills Almost in every organization, there is a need for SQL developer. Below are the skills that are in demand: 1. Back-end Developer Unlike front-end developer who managed the look and feels of a web app, back-end developer has to manage the data show to the user is proper and data updated in the database tables are correct. 2. Database Administrator A database administrator is someone who managed the database and its objects. DBA is the one who decides on access right of users. 3. Data Analyst The data analyst is the one who analyzes the data for a meaningful output. Why should we use SQL For almost every application data is important. To store and manage we need a database. And to access, use and manipulate that data we need a standard language. SQL is easy to learn a language, use to manage data stored in the database. One can learn the basics of SQL within a few days. It can be embedded in any other programming language. It is easy to code in SQL. Complex queries can be written in few lines of code. Hence SQL should be used for database related tasks. With the growing importance of data in the present era, the importance and need of someone who can understand and play with data are also increasing. SQL is getting extended to cloud platforms. Now one can query over millions and trillions of records in no time. It is also used in cutting edge technology like data science. Hence deep knowledge of SQL and it’s services can land you up in one the highest paying jobs. Who is the right audience for learning SQL technologies Anyone who has an interest in playing with data is the right audience for learning SQL technologies. Someone who enjoys analyzing data and getting something meaningful out of it. How SQL helps in career growth Learning SQL might help you land up in hot jobs like data scientist and data analyst. It opens door to cloud platforms as well. Database administrator and database architect are offered attractive pay scales by a reputed organization. SQL is an old but important language. It provides you the capability to store and manage data. It gives you all the powers to deal with relational data. It’s simple to learn but might get you an attractive job offer from reputed organization. This has been a guide to What is SQL. Here we discussed the various SQL subsets and top SQL companies with advantage and scope. You can also go through our other suggested articles-
What is the Kuiper Belt?Edit The Kuiper Belt (pronounced KYE-PER) is a disk of icy debris extending outwards from the orbit of Neptune at 30 AU to 55 AU from the Sun. Also known as the Edgeworth-Kuiper Belt, this region had been hypothesized variedly for decades but it was only in 1992 that direct evidence of its existence was found for the first time. However, there is still uncertainty who deserves credit for its discovery. Similar to the asteroid belt, the Kuiper Belt is mainly composed of small bodies. It is highly likely that the Kuiper Belt objects are primitive remnants from the formation of the solar system about 4.6 billion years ago. Frozen volatiles or "ices" such as water, ammonia, and methane are what comprise the Kuiper belt. It is believed that the belt and the even more distant Oort Cloud are the sources of comets that orbit the Sun , acting as a resevoir for these objects. Since 2006 there have been discoveries of several other kuiper belts in other solar systems dividing into two types: wide belts, with a radii of over 50 AU, and narrow belts, with a radii ranging from 20 AU to 30 AU. Discovery of the Kuiper BeltEdit - 1943: Astronomer Kenneth Edgeworth suggests that a reservoir of comets and larger bodies resides beyond the planets. - 1950: Astronomer Jan Oort theorizes that a vast population of comets may exist in a huge cloud on the distant edges of our solar system. - 1951: Astronomer Gerard Kuiper predicts the existence of a belt of icy objects just beyond the orbit of Neptune. - 1992: After five years of searching, astronomers David Jewitt and Jane Luu discover the first KBO (Kuiper Belt object), 1992QB1. KBOs - Kuiper Belt ObjectsEditA Kuiper Belt Object is any of many minor planets in the Kuiper belt outside the orbit of Neptune at the edge of the solar system. Since the belt was discovered in 1992, the number of known Kuiper belt objects (KBOs) has increased to over a thousand, and more than 70,000 KBOs are over 100 km (62 mi) in diameter. Some of the largest KBOs include Quaoar, Haumea, Eris, Pluto, and Makemake. Most of the large bodies of the Kuiper Belt consist mainly of ice and gases. Quaoar- Discoverd in 2002, Quaoar's diameter is roughly 1300 km (800 mi). Quaoar is about 4 billion miles away from Earth, which is the farthest object in the solar system to be resolved by a telescope. Haumea- Discovered in 2003, Haumea is believed to have a mass that is one-third that of Pluto and is roughly the same size as Pluto. It takes 285 years for Haumea to make one orbit around the Sun. Haumea is one of the fastest rotating objects in our solar systen and completes a full turn on its axis every 4 hours; Due to the fast rotation, the planet has become elongated in shape and is one of the strangest shaped dwarf planets that has been discovered in our Solar System.Eris- The detection of Eris in 2005 provoked debate about Pluto's classification as a planet. It has a diameter between 2400 and 3000 km (1490 mi to 1860 mi) and is 27% more massive than Pluto. Eris is the most distant object ever seen in orbit around the Sun. It is almost 10 billion miles from the Sun, more than 3 times more distant than Pluto, and takes more than twice as long to orbit the Sun as Pluto; It has an orbital period of 556.7 years. Pluto - In August 2006, Pluto was demoted to under the classification as a "dwarf planet." Discovered in 1930, Pluto is still considered one of the largest members of the Kuiper Belt region. The dwarf planet is about 1413 miles in diameter and 3.7 billion miles away from the Sun. Makemake- Currently, the second brightest Kuiper Belt object to be seen by the naked eye after Pluto is Makemake; Makemake's diameter is estimated to be between 1,360-1,480 km (smaller than Pluto and Eris). Because of its distance, Makemake takes about 310 years to orbit the Sun. Ninth Planet- On January 20, 2016 astronomers have found evidence, in terms of many calculations and simulations, revealing there to be a super massive planet hiding in the Kuiper belt. It has a mass 10 times that of our own Earth and is orbiting the sun in an extremely elongated elliptical that is 20 times farther away from the sun than Neptune is. It's existence would explain a few features of objects in the Kuiper Belt for example many objects in the Kuiper Belt are tilted to the same degree which rarely occurs at random. Also on the subject of any other objects colliding with the massive planet it seems to effectively control the other objects' orbits with it gravitational pull and slow orbit taking 10 to 20 thousand years. - Beyond the orbit of Neptune - An icy body Woosah (Project Plan) I plan on keeping some of the existing information but rephrase and reorder most of it so that it has some sense of organization. I will add more up-to-date pictures and insert at least one video. Most of the information is very general, so I'd like to go into detail about certain aspects of the belt. Besides minor spelling and grammatical errors that I want to fix, I'm going to clean up the layout of the page so that it is more organized. I am also going to try to make the descriptions and information easier to understand and follow so that it is clear and concise. I will also take out any unneccesary information that I feel strays from the main focus of the page.
Gravitational Waves May Have Come From a Black Hole Pair Born in a "Mosh Pit" Back in February, astronomers detected the world's first signal of gravitational waves — ripples through the fabric of space-time that come from powerful events like exploding stars. In the case of the signal detected in February, the gravitational waves came from a pair of black holes that spiraled closer and closer together until they merged. But how did those black holes end up so close together in the first place? New research published in the Astrophysical Journal Letters might have an answer. Dense pockets of stars called globular clusters are basically black hole factories. Stars inside them will regularly collapse into black holes. The really big black holes get sucked into the middle of the cluster, and they swirl around like ingredients in a blender. Pairs of black holes will sometimes get ejected during the shuffle. "Once those kick each other out, either as single or binary black holes, the next-most massive black holes rush in and start a new mosh pit," astronomer Carl Rodriguez told New Scientist. The ones that get ejected as pairs might be close enough to merge and produce gravitational waves. When Rodriguez and his team simulated how black holes behave inside globular clusters, they found that of the 262 black hole pairs that came out of their simulation, 14 of those were about the same mass as the two black holes that produced the gravitational waves in February. The research doesn't prove that the binary black holes that created the first gravitational wave signal were created this way. They could have also formed from a separate pair of stars that collapsed into black holes outside the cluster. If the stars were originally close enough together, they could have formed a binary black hole that would eventually merge. Astronomers will be able to study the source of gravitational waves once we've detected more of them. The Laser Interferometer Gravitational Wave Observatory is searching for more signals.
Presentation on theme: "Think Like a Scientist! Nature of Science"— Presentation transcript: 1 Think Like a Scientist! Nature of Science Big Idea 1: The Practice of ScienceBig Idea 2: The Characteristics of Scientific KnowledgeMary Tweedy, Curriculum Support SpecialistKeisha Kidd, Curriculum Support SpecialistDr. Millard Lightburn, Instructional Supervisor 2 Grade 5 Pacing Guide Topic 2: Thinking Like a Scientist SC.5.N.1.1 Define a problem, use appropriate reference materials to support scientific understanding, plan and carry out scientific investigations of various types such as: systematic observations, experiments requiring the identification of variables, collecting and organizing data, interpreting data in charts, tables, and graphics, analyze information, make predictions, and defend conclusions. AASC.5.N.2.1 Recognize and explain that science is grounded in empirical observations that are testable; explaining must always be linked with evidence. AASC.5.N.2.2 Recognize and explain that when scientific investigations are carried out, the evidence produced by those investigations should be replicable by others. AASC.5.P.8.1 Compare and contrast the basic properties of solids, liquids, and gases, such as mass, volume, color, texture, and temperature.MACC.5.MD.1.2 Make a line plot to display a data set of measurements in fractions of a unit (1/2, 1/4, 1/8). Use operations on fractions for this grade to solve problems involving information presented in line plots. For example, given different measurements of liquid in identical beakers, find the amount of liquid each beaker would contain if the total amount in all the beakers were redistributed equally.Department of Mathematics and Science 3 Engage: Have students sing the rap to identify key steps in the scientific method. Students share out. 4 THE SCIENTIFIC METHOD A process or steps scientists can use to gather information and answer questions!Engage/Explore: play the video clip. Ask students to define the scientific method in one sentence in their notebooks. Share out and discuss. Then click to reveal one way to define it. 5 Scientific Method Overview Make ObservationsAsk questionsDo Background ResearchForm a Hypothesis that is TestableExperiment to Test your HypothesisAnalyze Results & Draw ConclusionsAsk students to name the steps. 6 Scientific Method Study Jams – Scientific Method Video & Karaoke Song Learn how to think and solve problems like ascientist when Tim and Moby explore scientificmethods in this Brain POP movie.Hyperlinks to Online resources 7 Step 1: Ask a Question What do you want to find out? Identify one question thatcan be answered byperforming an experiment.An experiment is a set of steps you follow to test a hypothesis.This question will be the Problem Statement.Discuss 8 Step 2: Make a Hypothesis Look at the Problem Statement and identify the one factor that can be tested. This is the manipulated or independent variable.Form an idea or educated prediction that can be tested by an experiment.Write down your Hypothesis: “If (I do this) then (this) will happen.”Discuss 9 Step 3: Plan the Investigation Identify and record the factors that can affectthe results of the experiment under Variables.1. Test (independent/manipulated) variable or the factor that is changed in the experiment.(See previous Step 3.)2. Constant variables or all the factors to be kept the same (controlled) in the experiment.3. Outcome (dependent/responding) variable or the data to be collected during the experiment.4. Control Group (not found in all experiments) -A group that is untreated by the factor being testedthat serves as a reference for comparison to theexperimental group.Discuss 10 Step 4: Planning Continued Write your procedures or the steps you willfollow in your experiment.Each procedure step needs to be numbered.Each step needs to begin with a verb.These procedures will insure that all variables are kept the same (constant) or controlled except the one you are testing (independent).Identify control group = what remains the same (not the test variable).Figure out and collect the materials needed for the experiment.Discuss 11 Step 5: Collect, Organize, and Display Data Start the experiment.Observe and record the quantitative data (numbers or measurements) collected during the experiment on a data table. (evidence)Repeat the experiment three or more times to confirm results.Take pictures during the experiment.Graph your data from all trials. (Dry Mix)Display under Data.Restate your data in a narrative form under results.Discuss 12 Step 6: Drawing Conclusions What was investigated? (Describe the problem statement.)Restate your hypothesis, and tell if it was supported(true) or not supported (false).What were the major findings – the evidence ? (Explain your results.)Look at everything that may have affected your results. What possible explanation can you offer for your findings?Discuss 13 Step 7: Making Applications What recommendations do you have for further study and for improving the experiment?Explain what you learned from your experiment that could be applied in real life.List any new question(s) that your experiment lead you to ask that could be tested in a new investigation.Discuss 14 Mysterious M&M’s Take a closer look at an M&M 1. What are some things you observe about an M&M?Record your observations in your notebook.2. Talk in your group and discuss some of the properties youobserved about the M&M’s.3. Do you have both Qualitative and Quantitative Observations?4. Break open the M&M and look inside.5. Describe what you observe in words and make a drawing toshow what the inside of the M&M looks like.Explore: Distribute an M&M to each student. As students begin to explore the characteristics of an M&M, listen to the conversations that student groups are having. Students should identify properties such as the following:Size, shape, color, and textureDifferent colored layers on the insideExplain to students that their descriptions of M&M’s are all properties of M&M’s. Read the story from the Student Activity sheet, (see next slide) the student noticed that the color came off of an M&M when it fell in the water. Ask students if they ever had their M&M’s get wet and start to lose their color.You can click on the hyperlink: Adapted from Inquiry in Action Chapter 1 Activity 1.1 for the original lesson.Adapted from Inquiry in Action Chapter 1 Activity 1.1 15 An M&M ExperienceSometimes you can learn a lot about something by looking at it very closely or in ways you haven’t looked at it before. You may even discover things kind of by accident. This is what happened to me the other day when I was eating some M&M’s and drinking a cup of water. I was almost done when one of my M&M’s fell into the water that was left in my cup. I didn’t care too much because I could eat that one even though it was wet. I decided to eat it but when I began to reach into the cup to take it out, I was kind of surprised by what I saw. There was an area of color in the water around the M&M, which I guess had dissolved into the water.Explore continued: Tell students that in the Try this activity, they will see what happens to the sugar and color coating of an M&M when it is placed in a plate of water. Have students place an M&M in a dish of water and observe.Giving students an opportunity to observe an M&M in water will give them the context and motivation to want to find out more about how M&M colors look when they dissolve in water. From this experience, you can get them to ask questions that they can investigate. Students will conduct the following procedure and record their observations.ProcedurePour enough room-temperature water into a white plastic or foam plate so that the water is deep enough to completely cover an M&M.Once the water has settled, place 1 M&M in the center of the plate. Be careful to keep the water and M&M as still as possible. Observe for about 1 minute. 16 Mysterious M&M’sLet’s try placing an M&M in water to get a better ideaof what the student in the story observed.Read the procedures on the lab sheet.What materials will be needed for each group ?- container of room temperature water- one M&M- white foam plateFollow the procedures #’sDiscuss your group’s observations.What do you notice about the movement of thecolor from the M&M?Explore/Explain: Have students compare their results.Ask students: what do they notice about the movement of the color from their M&M? What questions do they have?Expected results: Each colored coating of M&M will dissolve in a circular pattern around the M&M. Students may also mention the white streaks in the water from the sugar coating. If anyone notices differences such as “the color moved over to one side more than the other,” check to see that the plate is level.Point out to students that because the water makes the colored coating come off the M&M and mix into the water, the water is dissolving the sugar and color. Because the colored coating on M&M’s dissolves in a similar pattern each time one is placed in water, this is a characteristic property of the M&M coating.Empty the plate of water and M&M into a bucket, bowl, or sink. Dry the plate with a paper towel. 17 M&M Questions to Investigate Look at the variables below:color of M&M’snumber of M&M’stemperature of watertype of liquidWhat question(s) can be investigated by changingone variable above at a time?Engage/Explore: Tell students these variables can be changed to learn more about M&M’s. Have them work as a team and make a list of questions that could be investigated. Write questions on the board or chart paper and discuss. Tell students they will investigate one of the questions. 18 Question to Investigate through an Experiment (Problem Statement)Does the temperature of the wateraffect how fast the coloredcoating dissolves from an M&M?Explore: Tell students that this is the first question that will be investigated. It will be investigated by the class with each group doing one trial. Each group will share their data collected with the whole class. The investigation will be done in the format of doing a science fair project.You can click on the hyperlink: Adapted from Inquiry in Action Chapter 1 Activity 1.5 for the original lesson.Adapted from Inquiry in Action Chapter 1 Activity 1.5 19 Hypothesis Write your own hypothesis. If ………………………, then……………….. Explore continued: Sample hypothesis: If an M&M is placed in three different temperatures of water, room temperature, cold, and hot water, then the M&M placed in hot water will dissolve the fastest. 20 Materials 3 Same-color M&M’s 1 White foam dessert plates (with measurements)Room-temperature waterHot waterCold waterMeasuring cupCentimeter rulerLarge container or bowlPaper towelsGraduated cylinderStopwatch(stopwatch online)Explore continued: Have each group member assigned to a lab role: Project Director (PD), Materials Manager (MM), Technical Manager (TM), and a Safety Director (SD) or a variation of these roles.Materials Manager will collect the materials and getting the type of water when needed.Each group needs these materials. You can have a thermos of hot water, a pitcher of cold and room temperature water set up in one location for students use.Students can use a stopwatch or timer or go online for a stopwatch. 21 Procedures Take the temperature of the room temperature water. Pour 50mL of room-temperature water into the plate.Place a same-colored M&M in the center of the plate with the help of your partners, and observe for 1 minute. (stopwatch online)Record the qualitative measurements you see in the 1 min.Measure the distance in centimeters that the colored coating traveled. (measure from the center).Record your quantitative observations on the group Data Table Trial #___.Repeat steps 1-6 with the cold water, then with hot water.Record your group’s trial data on the class data chart.Copy the other groups’ trial data.Find the average for all of the trials and record.Explore/Explain: Students can use a stopwatch or timer or go online for a stopwatch. 22 Distance Coating Traveled in Centimeters Class Data Collection:Distance Coating Traveled in CentimetersTrialRoom TemperatureWaterColdHot#1cm#2#3#4#5#6AverageExplain: Have each group create a graph to display the data. 23 Share Each Group’s Trial Data Explain/Evaluate: Have students share their data. Ask students whether they noticed a difference in the movement of color in the different temperatures of water. Discuss with students how to write the results on the lab sheet. 24 Conclusion What was investigated? (Describe the problem statement.) Restate your hypothesis, and tell if it was supported or not supported.What were the major findings? (Explain your evidence.)Explain/Evaluate: Students answer the questions to write a conclusion. 25 ApplicationIf the experiment was to be repeated should anything be done differently?Explain what you learned from your experiment that could be applied in a real life situation.List any new questions that your experiment lead you to ask that could be tested in a new investigation.Explain/Evaluate: Students answer the questions to write an application. 26 Communicate Your Inquiry You can use the science fair blank template to create a Power Point presentation.You can duplicate your Power Point presentation and display on a mini-Science Fair Project Board.Explain/Evaluate: Show students the Science Fair PP template and how to use it. 27 New M&M Investigation?Form a new question or state a New Problem on the same topic.What do you still want to know?What more can you learn?Extend: Students can do a new M&M investigation. 28 Scientific Method Review Use as needed for review and also for Science Fair Background preparation. 29 1. Ask a Question or State a Problem Asking WHAT? or HOW? about something you observedReview 30 2. Research your TopicGather information that will help you answer your question.Library, Internet, Interviews, ExperimentsReview 31 An educated guess! 3. State your HYPOTHESIS A Hypothesis is an explanation for a question that can be formally tested.An educated guess!If…then…Review 32 4. Design an ExperimentA procedure is a set of directions designed to test your Hypothesis… Is it is true or false.A procedure must be repeatable, and easy to understand for others to duplicate.Each procedure step needs to be numbered.Each step needs to begin with a verb.Figure out and collect the materials needed for the experiment.Review 33 Step 4 Planning Continued Write your procedures or the steps you willfollow in your experiment.Each procedure step needs to be numbered.Each step needs to begin with a verb.These procedures will insure that all variables are kept the same (constant) or controlled except the one you are testing.Figure out and collect the materials needed for the experiment.Review 34 In a well designed Experiment, you need to keep all variables the same except one. Test/Independent/Manipulated Variable: (CAUSE)The factor that is changed in an experiment…it is what you are testing!Constant/Controlled Variable(s):The factor(s) that remains the same!Outcome/Dependent/Responding Variable: (EFFECT) The data you collectReview 35 5. Conduct your Experiment Perform your experiment by following your written procedure.Be sure to follow all safety rules!Review 36 6. Collect DataThe observations and measurements you make in an experiment are called Data.Review 37 7. Analyze Data Did your experiment support your hypothesis? What happened during your experiment?Does additional research need to be conducted?Review 38 8. Conclusion Does your data and observations support your hypothesis? “My hypothesis was(supported or not supported)because __________”Review 39 Share your results and data with others. 9. CommunicationShare your results and data with others.Sources: written, spoken, video, TV, papers, lecture . . .Review 40 Communicate Your Inquiry You can use the science fair blank template to create a Power Point presentation.You can duplicate your Power Point presentation and display on a Science Fair Project Board.ReviewClick on hyperlink: science fair blank template to show students a resource they can modify to create their own science investigation presentation. 41 10. New ProblemForm a new question or state a New Problem on the same topic.What more can you learn?What do you still want to know?Review 42 The Science Fair is Coming! Think Like a Scientist!Ask questions InvestigateObserve ExperimentThe Science Fair is Coming! 43 Scientific Method - Resource Links in Action Chapter 1 Molecules in Motion : Activities 1.1 and 1.5)Experiment VocabularyControlExperimentThinking Like a ScientistFundamentals of Experimental DesignResources on line
You may have heard the saying "You can prove anything with statistics," which implies that statistical analysis cannot to be trusted, that the conclusions that can be drawn from it are so vague and ambiguous that they are meaningless. Yet the opposite is also true. Statistical analysis can be reliable and the results of statistical analysis can be trusted if the proper conditions are established. What Is Statistical Analysis? Statistical analysis uses inductive reasoning and the mathematical principles of probability to assess the reliability of a particular experimental test. Mathematical techniques have been devised to allow measurement of the reliability (or fallibility) of the estimate to be determined from the data (the sample, or "N") without reference to the original population. This is important because researchers typically do not have access to information about the whole population, and a sample—a subset of the population— is used. Statistical analysis uses a sample drawn from a larger population to make inferences about the larger population. A population is a well-defined group of individuals or observations of any size having a unique quality or characteristic. Examples of populations include first-grade teachers in Texas, jewelers in New York, nurses at a hospital, high school principals, Democrats, and people who go to dentists. Corn plants in a particular field and automobiles produced by a plant on Monday are also populations. A sample is the group of individuals or items selected from a particular population. A random sample is taken in such a way that every individual in the population has an equal opportunity to be chosen. A random sample is also known as an unbiased sample. Most mail surveys, mall surveys, political telephone polls, and other similar data gathering techniques generally do not meet the proper conditions for a random, unbiased sample, so their results cannot to be trusted. These are "self-selected" samples because the subjects choose whether to participate in the survey and the subjects may be picked based on the ease of their availability (for example, whoever answers the phone and agrees to the interview). Selecting a Random Sampling The most important criterion for trustworthy statistical analysis is correctly choosing a random sample. For example, suppose you have a bucket full of 10,000 marbles and you want to know how many of the marbles are red. You could count all of the marbles, but that would take a long time. So, you stir the marbles thoroughly and, without looking, pull out 100 marbles. Now you count the red marbles in your random sample. There are 10. Thus you could conclude that approximately 10 percent of the original marbles are red. This is a trustworthy conclusion, but it is not likely to be exactly right. You could improve your accuracy by counting a larger sample; say 1,000 marbles. Of course if you counted all the marbles, you would know the exact percentage, but the point is to pick a sample that is large enough (for example, 100 or 1,000) that gives you an answer accurate enough for your purposes. Suppose the 100 marbles you pulled out of the bucket were all red. Would this be proof that all 10,000 marbles in the bucket were red? In science, statistical analysis is used to test a hypothesis . In the example we are testing, the hypothesis would be "all the marbles in the bucket are red." Statistical inference makes it possible for us to state, given a sample size (100) and a population size (10,000), how often false hypotheses will be accepted and how often true hypotheses are rejected. Statistical analysis cannot conclusively tell us whether a hypothesis is true; only the examination of the entire population can do that. So "statistical proof" is a statement of just how often we will get "the right answer." Using Basic Statistical Concepts Statistics is applicable to all fields of human endeavor, from economics to education, politics to psychology. Procedures worked out for one field are generally applicable to the other fields. Some statistical procedures are used more often in some fields than in others. Example 1. Suppose the Wearemout Pants Company wants to know the average height of adult American men, an important piece of information for a clothing manufacturer producing pants. The population is all men over the age of 25 who live in the United States. It is logistically impossible to measure the height of every man who lives in the United States, so a random sample of around 1,000 men is chosen. If the sample is correctly chosen, all ethnic groups, geographic regions, and socioeconomic classes will be adequately represented. The individual heights of these 1,000 men are then measured. An average height is calculated by dividing the sum of these individual heights by the total number of subjects (N = 1,000). By doing so, imagine that we calculate an average height is 1.95 meters (m) for this sample of adult males in the United States. If a representative sample was selected, then this figure can be generalized to the larger population. The random sample of 1,000 men probably included some very short men and some very tall men. The difference between the shortest and the tallest is known as the "range" of the data. Range is one measure of the "dispersion" of a group of observations. A better measure of dispersion is the "standard deviation." The standard deviation is the square root of the sum of the squares of the differences divided by one less than the number of observations. In this equation, xi is an observed value and is the arithmetic mean. In our example, if a smaller height interval is used (1.10 m, 1.11 m, 1.12 m, 1.13 m, and so on) and the number of men in each height interval plotted as a function of height a smooth curve can be drawn which would have a characteristic shape, known as a "bell" curve or "normal frequency distribution." A normal frequency distribution can be stated mathematically as The value of sigma (σ) is a measure of how "wide" the distribution is. Not all samples will have a normal distribution, but many do, and these distributions are of special interest. The following figure shows three normal probability distributions. Because there is no skew, the mean, median, and mode are the same. The mean of curve (a) is less than the mean of curve (b), which in turn is less than the mean of (c). Yet the standard deviation, or spread, of (c) is least, whereas that of (a) is greatest. This is just one illustration of how the parameters of distributions can vary. Example 2. One of the most common uses of statistical analysis is in determining whether a certain treatment is efficacious. For example, medical researchers may want to know if a particular medicine is effective at treating the type of pain resulting from extraction of third molars (known as "wisdom" teeth). Two random samples of approximately equal size would be selected. One group would receive the pain medication while the other group received a "placebo," a pill that looked identical but contained only inactive ingredients. The study would need to be a "double-blind" experiment, which is designed so that neither the recipients nor the persons dispensing the pills knew which was which. Researchers would know who had received the active medicines only after all the results were collected. Example 3. Suppose a student, as part of a science fair project, wishes to determine if a particular chemical compound (Chemical X) can accelerate the growth of tomato plants. In this sort of experiment design, the hypothesis is usually stated as a null hypothesis : "Chemical X has no effect on the growth rate of tomato plants." In this case, the student would reject the null hypothesis if she found a significant difference. It may seem odd, but that is the way most of the statistical tests are set up. In this case, the independent variable is the presence of the chemical and the dependent variable is the height of the plant. The next step is experiment design. The student decides to use height as the single measure of plant growth. She purchases 100 individual tomato plants of the same variety and randomly assigns them to 2 groups of 50 each. Thus the population is all tomato plants of this particular type and the sample is the 100 plants she has purchased. They are planted in identical containers, using the same kind of potting soil and placed so they will receive the same amount of light and air at the same temperature. In other words, the experimenter tries to "control" all of the variables, except the single variable of interest. One group will be watered with water containing a small amount of the chemical while the other will receive plain water. To make the experiment double-blind, she has another student prepare the watering cans each day, so that she will not know until after the experiment is complete which group was receiving the treatment. After 6 weeks, she plans to measure the height of the plants. The next step is data collection. The student measures the height of each plant and records the results in data tables. She determines that the control group (which received plain water) had an average (arithmetic mean) height of 1.3 m (meters), while the treatment group had an average height of 1.4 m. Now the student must somehow determine if this small difference was significant or if the amount of variation measured would be expected under no treatment conditions. In other words, what is the probability that 2 groups of 50 tomato plants each, grown under identical conditions would show a height difference of 0.1 m after 6 weeks of growth? If this probability is less than or equal to a certain predetermined value, then the null hypothesis is rejected. Two commonly used values of probability are 0.05 or 0.01. However, these are completely arbitrary choices determined mostly by the widespread use of previously calculated tables for each value. Modern computer analysis techniques allow the selection of any value of probability. The simplest test of significance is to determine how "wide" the distribution of heights is for each group. If there is a wide variance (σ) in heights (say, σ = 25), then small differences in mean are not likely to be significant. On the other hand, if the dispersion is narrow (for example, if all the plants in each group were close to the same height, so that σ = 0.1) then the difference would probably be significant. There are several different tests the student could use. Selecting the right test is often a tricky problem. In this case, the student can reject several tests outright. For example, the chi-square test is suitable for nominal scales (yes or no answers are one example of nominal scales), so it does not work here. The F -test measures variability or dispersion within a single sample. It too is not suitable for comparing two samples. Other statistical tests can also be rejected as inappropriate for various reasons. In this case, since the student is interested in comparing means, the best choice is a t test. The t -test compares two means using this formula: In this case, the null hypothesis assumes that μ1 − μ2 = 0 (no difference in the sample groups), so that we can say: The quantity is known as the standard error of the mean difference. When the sample sizes are the same, . The standard error of the mean difference is the square root of the sums of the squares of the standard errors of the means for each group. The standard error of the mean for each group is easily calculated from . N is the sample size, 50, and the student can calculate the standard deviation by the formula for standard deviation given above. The final experimental step is to determine sensitivity. Generally speaking, the larger the sample, the more sensitive the experiment is. The choice of 50 tomato plants for each group implies a high degree of sensitivity. Students, teachers, psychologists, economists, politicians, educational researchers, medical researchers, biologists, coaches, doctors and many others use statistics and statistical analysis every day to help them make decisions. To make trustworthy and valid decisions based on statistical information, it is necessary to: be sure the sample is representative of the population; understand the assumptions of the procedure and use the correct procedure; use the best measurements available; keep clear what is being looked for; and to avoid statements of causal relations if they are not justified. see also Central Tendency, Measures of; Data Collection and Interpretation; Graphs; Mass Media, Mathematics and the. Huff, Darrell. How to Lie With Statistics. New York: W. W. Norton & Company, 1954. Kirk, Roger E. Experimental Design, Procedures for the Behavioral Sciences. Monterrey, CA: Brooks/Cole Publishing Company, 1982. Paulos, John Allen. Innumeracy: Mathematical Illiteracy and Its Consequences. New York: Hill & Wang, 1988. Tufte, Edward R. The Visual Display of Quantitative Information. Cheshire, CT: Graphics Press, 1983. MEAN, MEDIAN, AND MODE The average in the clothing example is known as the "arithmetic mean," which is one of the measures of central tendency. The other two measures of central tendency are the"median" and the "mode." The median is the number that falls in the mid-point of an ordered data set, while the mode is the most frequently occurring value. "Statistical Analysis." Mathematics. . Encyclopedia.com. (December 12, 2017). http://www.encyclopedia.com/education/news-wires-white-papers-and-books/statistical-analysis "Statistical Analysis." Mathematics. . Retrieved December 12, 2017 from Encyclopedia.com: http://www.encyclopedia.com/education/news-wires-white-papers-and-books/statistical-analysis Classical Statistical Analysis Classical Statistical Analysis Classical statistical analysis seeks to describe the distribution of a measurable property (descriptive statistics) and to determine the reliability of a sample drawn from a population (inferential statistics). Classical statistical analysis is based on repeatedly measuring properties of objects and aims at predicting the frequency with which certain results will occur when the measuring operation is repeated at random or stochastically. Properties can be measured repeatedly of the same object or only once per object. However, in the latter case, one must measure a number of sufficiently similar objects. Typical examples are measuring the outcome of tossing a coin or rolling a die repeatedly and count the occurrences of the possible outcomes as well as measuring the chemical composition of the next hundred or thousand pills produced in the production line of a pharmaceutical plant. In the former case the same object (one and the same die cast) is “measured” several times (with respect to the question which number it shows); in the latter case many distinguishable, but similar objects are measured with respect to their composition which in the case of pills is expected to be more or less identical, such that the repetition is not with the same object, but with the next available similar object. One of the central concepts of classical statistical analysis is to determine the empirical frequency distribution that yields the absolute or relative frequency of the occurrence of each of the possible results of the repeated measurement of a property of an object or a class of objects when only a finite number of different outcomes is possible (discrete case). If one thinks of an infinitely repeated and arbitrarily precise measurement where every outcome is (or can be) different (as would be the case if the range of the property is the set of real numbers), then the relative frequency of a single outcome would not be very instructive; instead one uses the distribution function in this (continuous) case which, for every numerical value x of the measured property, yields the absolute or relative frequency of the occurrence of all values smaller than x. This function is usually noted as F (x ), and its derivative F’ (x ) = f (x ) is called frequency density function. If one wants to describe an empirical distribution, the complete function table is seldom instructive. This is why the empirical frequency or distribution functions are often represented by a few parameters that describe the essential features of the distribution. The so-called moments of the distribution represent the distribution completely, and the lower-order moments represent the distribution at least in a satisfactory manner. Moments are defined as follows: where k is the order of the moment, n is the number of repetitions or objects measured, and c is a constant that is usually either 0 (moment about the origin) or the arithmetic mean (moment about the mean), the first-order mean about the origin being the arithmetic mean. In the frequentist interpretation of probability, frequency can be seen as the realization of the concept of probability: It is quite intuitive to believe that if the probability of a certain outcome is some number between 0 and 1, then the expected relative frequency of this outcome would be the same number, at least in the long run. From this, one of the concepts of probability is derived, yielding probability distribution and density functions as models for their empirical correlates. These functions are usually also noted as f (x ) and F (x ), respectively, and their moments are also defined much like in the above formula, but with a difference that takes into account that there is no finite number n of measurement repetitions: where the first equation can be applied to discrete numerical variables (e.g., the results of counting), while the second equation can be applied to continuous variables. Again, the first-order moment about 0 is the mean, and the other moments are usually calculated about this mean. In many important cases one would be satisfied to know the mean (as an indicator for the central tendency of the distribution) and the second-order moment about the mean, namely the variance (as the most prominent indicator for the variation). For the important case of the normal or Gaussian distribution, these two parameters are sufficient to describe the distribution completely. If one models an empirical distribution with a theoretical distribution (any non-negative function for which the zero-order moment evaluates to 1, as this is the probability for the variable to have any arbitrary value within its domain), one can estimate its parameters from the moments of the empirical distributions calculated from the finite number of repeated measurements taken in a sample, especially in the case where the normal distribution is a satisfactory model of the empirical distribution, as in this case mean and variance allow the calculation of all interesting values of the probability density function f (x ) and of the distribution function F (x ). Empirical and theoretical distributions need not be restricted to the case of a single property or variable, they are also defined for the multivariate case. Given that empirical moments can always be calculated from the measurements taken in a sample, these moments are also results of a random process, just like the original measurements. In this respect, the mean, variance, correlation coefficient or any other statistical parameter calculated from the finite number of objects in a sample is also the outcome of a random experiment (measurement taken from a randomly selected set of objects instead of exactly one object). And for these derived measurements theoretical distributions are also available, and these models of the empirical moments allow the estimation with which probability one could expect the respective parameter to fall into a specified interval in the next sample to be taken. If, for instance, one has a sample of 1,000 interviewees of whom 520 answered they were going to vote for party A in the upcoming election, and 480 announced they were going to vote for party B, then the parameter πA—the proportion of A-voters in the overall population—could be estimated to be 0.52, but this estimate would be a stochastic variable, which approximately obeys a normal distribution with mean 0.52 and variance 0.0002496 (or standard deviation 0.0158), and from this result one can conclude that another sample of another 1,000 interviewees from the same overall population would lead to another estimate whose value would lie within the interval [0.489, 0.551] (between 0.52 ± 1.96 0.0158) with a probability of 95 percent (the so-called 95 percent confidence interval, which in the case of the normal distribution is centered about the mean with a width of 3.92 standard deviations). Or, to put it in other words, the probability of finding more than 551 A-voters in another sample of 1,000 interviewees from the same population is 0.025. Bayesian statistics, as opposed to classical statistics, would argue from the same numbers that the probability is 0.95 that the population parameter falls within the interval [0.489, 0.551]. SEE ALSO Bayesian Statistics; Descriptive Statistics; Inference, Bayesian; Inference, Statistical; Sampling; Variables, Random Hoel, Paul G. 1984. Introduction to Mathematical Statistics. 5th ed. Hoboken, NJ: Wiley. Iversen, Gudmund. 1984. Bayesian Statistical Inference. Beverly Hills, CA: Sage. Klaus G. Troitzsch "Classical Statistical Analysis." International Encyclopedia of the Social Sciences. . Encyclopedia.com. (December 12, 2017). http://www.encyclopedia.com/social-sciences/applied-and-social-sciences-magazines/classical-statistical-analysis "Classical Statistical Analysis." International Encyclopedia of the Social Sciences. . Retrieved December 12, 2017 from Encyclopedia.com: http://www.encyclopedia.com/social-sciences/applied-and-social-sciences-magazines/classical-statistical-analysis "statistical analysis." A Dictionary of Computing. . Encyclopedia.com. (December 12, 2017). http://www.encyclopedia.com/computing/dictionaries-thesauruses-pictures-and-press-releases/statistical-analysis "statistical analysis." A Dictionary of Computing. . Retrieved December 12, 2017 from Encyclopedia.com: http://www.encyclopedia.com/computing/dictionaries-thesauruses-pictures-and-press-releases/statistical-analysis
Addition of vectors Addition of vectors-the two vectors a and b can be added giving the sum to be a + b this requires joining them head to tail we can translate the vector b till its tail meets the head of a the line segment that is directed from the tail of vector a to the head of vector b is the vector “a + b. In adding the east-west components of all the individual vectors, one must consider that an eastward component and a westward component would add together as a positive and a negative some students prefer to think of this as subtraction as opposed to addition. Vectors in standard position have a common origin and are used in the parallelogram rule of vector addition construct a parallelogram using two vectors in standard position the resultant is the diagonal of the parallelogram coming out of the common vertex. Addition of vectors sep 20, 2011 #1 asad1111 how to add two vectors knowing only angle between them for example if we have two vectors a and b having angle 40 degrees what will be the resultant sum of two this is the first question and he asked about 40 degree, this don't spell anything about 3 phase voltage. Adding and subtracting vectors to add or subtract two vectors, add or subtract the corresponding components let u → = 〈 u 1 , u 2 〉 and v → = 〈 v 1 , v 2 〉 be two vectors then, the sum of u → and v → is the vector. (a) addition or composition of vectors means finding the resultant of a number of vectors acting on a body (b) the vectors can be added geometrically and not algebraically (c) vectors, whose resultant is to be calculated behave independent of each other. Addition, subtraction and scalar multiplication of vectors, examples vectors - introduction there are physical quantities like force, velocity, acceleration and others that are not fully determined by their. C = plus(a,b) is an alternate way to execute a + b, but is rarely used it enables operator overloading for classes it enables operator overloading for classes examples. Example: given that , find the sum of the vectors solution: triangle law of vector addition in vector addition, the intermediate letters must be the same since pqr forms a triangle, the rule is also called the triangle law of vector addition graphically we add vectors with a head to tail approach. Stack exchange network consists of 174 q&a communities including stack overflow, the largest, most trusted online community for developers to learn, share their knowledge, and build their careers visit stack exchange. Triangle law of vector addition states that when two vectors are represented by two sides of a triangle in magnitude and direction taken in same order then third side of that triangle represents in magnitude and direction the resultant of the vectors. Draw the vectors in pictorial form choose any of the two vectors as the first and place the second vector’s tail on the head of the first draw a third vector joining the tail of the first to the head of the second this vector is called the “resultant” which represents the sum of the two. Advanced math solutions – vector calculator, simple vector arithmetic vectors are used to represent anything that has a direction and magnitude, length the most popular example of. Vectors and vector addition: a scalar is a quantity like mass or temperature that only has a magnitude on the other had, a vector is a mathematical object that has magnitude and direction a line of given length and pointing along a given direction, such as an arrow, is the typical representation of a vector. Learn how to add and subtract vectors by looking at free maths videos and example questions study the free resources during your math revision and pass your next math exam good luck and have fun. The graphical method for vector addition and scalar multiplication graphical addition consider the vectors u = (3, 4) and v = (4, 1) in the plane from the component method of vector addition we know that the sum of these two vectors is u + v = (7, 5)graphically, we see that this is the same as the result we would get by picking up one of the vectors (without changing either its direction. Addition of vectors Vector addition is the operation of adding two or more vectors together into a vector sum the so-called parallelogram law gives the rule for vector addition of two or more vectors for two vectors and , the vector sum is obtained by placing them head to tail and drawing the vector from the free tail to the free head. To add two vectors, a and b, we first break each vector into its components, ax and ay, and bx and by, as shown on the figure from the rules which govern the equality of vectors, the blue vector b is equal to the black vector b because it has equal equal length and equal direction. Finding the components of vectors for vector addition involves forming a right triangle from each vector and using the standard triangle trigonometry the vector sum can be found by combining these components and converting to polar form. Consider two vectors p and q acting on a body and represented both in magnitude and direction by sides oa and ab respectively of a triangle oablet θ be the angle between p and qlet r be the resultant of vectors p and qthen, according to triangle law of vector addition, side ob represents the resultant of p and q so, we have. - Vector addition: numerical calculates the magnitude and direction of the resultant given the magnitudes and directions of an arbitrary number of vectors to be added prerequisites students should understand the vector properties of magnitude and direction and be familiar with adding vectors graphically by the tip-to-tail method. - Vectors is the basic topic which helps us to make other topics in physics easier kinematics has a special topic named vector reversal method which is the easiest way of solving those sums. - Notes on vector addition name:_____ purpose: today i will learn vector components and how to add vectors together relevance: vectors are used in almost every chapter of physics learning to properly add vectors will make many problems much easier for you to solve success: i will be successful today if i can add horizontal, perpendicular, and odd angle vectors. Graphic addition permits one to see why mental addition is possible and why algebraic addition through components works so nicely to add vectors graphically, make a head-to-tail trail, then draw a short-cut arrow to connect the start to the finish. Experiment 2: vector addition uploaded by kamylle consebido this talks about the different methods of vector addition such as the polygon method, parallelogram method, and component method. In the above example, we demonstrated adding vectors physically by drawing and measuring them in the real world, we need much greater accuracy (however, you learned the important concept of visualizing vector addition) adding vectors mathematically. Vectors can be multiplied by a scalar to produce another vector multiplying vector x by 3 will give a new vector 3 times the length and parallel to x vector addition and subtraction.
The Moon. Photo: mmirrorless/Flickr, CC BY 2.0. When dabbling with the laws of motion in the 17th century, Isaac Newton realised that it’s possible to send an object from Earth into space. That is, he figured that if an object is shot away from Earth with enough velocity, it will reach space and start orbiting our planet. The Soviet Union achieved exactly this when they launched the Sputnik satellite in 1957 onboard a powerful rocket. For the first time, something on Earth had intentionally sent something else into space. Many of the satellites we launch today have little engines of their own, and by firing them1, they can escape Earth’s gravitational hold completely. That’s how we have sent missions to Venus, Mars, Saturn, even Pluto. The deceptively simple math that makes these feats of escape possible is captured in the Tsiolkovsky rocket equation, named for Soviet rocket scientist Konstantin Tsiolkovsky. This equation allows scientists and engineers to quantify and compare the energy required to reach different destinations in space. Its implications are far-reaching — but they are not entirely intuitive. Let’s see why. Getting to space The rocket equation tells us that the amount of energy a rocket must expend to go from Earth’s surface to low-Earth orbit 2 above is almost thrice as much as the energy needed to go from there to the Moon. Likewise, getting to this orbit costs more than twice the energy required to reach Mars from that orbit. (Energy expenditure is not the same as change in velocity, only proportional. But for simplicity’s sake, one has been used as an approximation of the other.) In fact, going from the ground to low-Earth orbit is 50% more expensive than starting in the low-Earth orbit and landing on the Moon. So lifting off from Earth and entering orbit is the biggest energy barrier to space exploration. In every rocket, the satellite it carries makes up a tiny fraction of the total rocket’s mass. However, the satellite’s small contribution makes a big difference. This is the core problem of rocket science. If you had to make a satellite more useful, you’d probably add more or bigger components to it, making it heavier. This in turn would mean the rocket has to carry a higher load, which then means it needs to carry more fuel. However, when more fuel is added, the total rocket becomes heavier, and then requires more fuel. And so on. As a thumb rule, the fuel requirement increases exponentially with every unit increase in the satellite’s mass. This is how we end up with rockets being mostly fuel, some metal and a small payload. Even the Saturn V rocket that sent astronauts to the Moon for the first time was 85% fuel, 13% rocket – including the chassis, its plumbing and parts. Only the remaining 2% was the Moonbound spacecraft and the astronauts inside. A bigger planet Now, say Earth is more massive than it currently is. This means it will exert a stronger gravitational pull on the rocket, so the rocket would have to expend more energy to leave for space. Which then means it will have to carry much, much more fuel than before. This way, if you keep increasing Earth’s mass, the fuel-to-mass ratio would start skyrocketing to a point where it would be simply impossible to build a functional rocket. A back of the envelope calculation suggests that if Earth was 50% more massive, you simply wouldn’t be able to get to space even with the most energetic fuel combination available in chemical rockets: liquid hydrogen and liquid oxygen (used in cryogenic engines). We know such planets already exist. Of the 4,000+ planets around other stars we have discovered to date in the Milky Way galaxy, about a thousand are super-Earths. These planets are up to 10-times more massive than Earth and up to twice as large. (Beyond these limits, planets can’t stay rocky and become gassy, like Jupiter and Neptune). Many of these super-Earths lie in the respective habitable zones around their stars, where conditions on the planets’ surfaces could support life as we know it. Since we have only searched a small fraction of our galaxy for planets, it’s fair to say there could be millions of super-Earths in the galaxy, many of which could host life. If intelligent life were to develop on these super-Earths, they would have a hard time building rockets that can take objects into space. And since even the most powerful cryogenic engines won’t get them to space, they might have to build something that can produce more thrust, like rockets propelled by nuclear reactions. These will likely be far more expensive than chemical rockets (assuming their economy is structured like Earth’s), but sometimes nature doesn’t give you a choice. Just like the rocket equation makes it exponentially harder to get off a planet, it also makes getting off objects with lesser gravity more easy. To understand how, let’s go to the Moon. Also read: The Ambitious Plan to Power Spacecraft to the Stars With Nuclear Reactions If humans wish to erect a permanent human settlement on Mars, say to improve the prospects of our survival in the super-long-term, we’ll need to send thousands of tonnes of material from Earth to the Martian surface on hundreds of rockets. This is what SpaceX hopes to do with its Starship vessel. The Moon’s gravity is much weaker than that of Earth. Moreover, the Moon lies at the outer edge of Earth’s gravity well, which means a rocket taking off from the Moon doesn’t have to work against Earth’s gravity as well. As a result, rockets need almost 5x less energy to lift off from the Moon’s surface than from Earth’s. So if we establish a settlement on the Moon first, we can eventually tap its resources to launch rockets from the Moon itself. NASA and ISRO missions have already discovered plenty of water ice near the Moon’s poles. Future human habitats built by mining the metal-rich lunar soil could tap into this water ice for their needs. The simple electrolysis technique could also split this water into hydrogen and oxygen, to be used separately to fuel cryogenic engines. Rockets lifting off from such an industrially enabled Moonbase can ride the lunar interplanetary highway to reach Mars more efficiently than from Earth. The capital cost would be enormously high, of course, but if governments have resolved to expand sustainably into the Solar System, they would also know that doing so will take many centuries, if not millennia. In this extended scheme of things, launching missions from the Moon instead of Earth will reduce the energy, and therefore resources, we will need. The Moon’s accessibility, low gravity-barrier and resource potential are the reasons why some argue that it’s important to return to the Moon before getting to Mars. The belt and beyond The Moon’s advantages may extend to making homes for ourselves in the outer Solar System as well. The rocket equation tells us that even if objects in the outer reaches of the system may be closer to Mars than to the Moon, the red planet has a deeper gravity well than the Moon, keeping the Moon advantageous. Specifically, a rocket going from Mars to the asteroid belt will need 40% more energy than when going from the Moon to the asteroid belt – even though Mars is about 75 million km closer to the belt. This is the difference gravity makes, and which the rocket equation allows us to see. The Moon can accelerate the expansion of settlements to encompass resource-rich asteroids. The bigger ones among them, like Ceres and Vesta, can in turn play the same role as the Moon can for Mars, and help humans move on to the moons of Jupiter and Saturn, and beyond. The most important takeaway from the rocket equation is that the ability to extract and harness raw materials from low-gravity, resourceful space objects frees us from having to drag everything out of Earth’s gravitational pull first. Ultimately, we can’t hope to journey between the stars if we don’t journey between the planets first, in an Earth-independent way. The author thanks Adithya K. Pani for reviewing the article. This article was originally published on Jatan Mehta‘s blog and has been republished here with permission.
This image shows galaxies clumped together in the Fornax cluster, located 60 million light-years from Earth. The picture was taken by WISE, but has been artistically enhanced to illustrate the idea that clumped galaxies will, on average, be surrounded by larger halos of dark matter (represented in purple). Image released May 22, 2014. Conventional thinking suggests that the most massive black holes possess a ringed doughnut-shaped torus of gas and dust trapped in orbit around them. But if we know one thing about black holes, they're anything but conventional. Now, astronomers have analyzed data from NASA's Wide-field Infrared Survey Explorer (WISE) of thousands of supermassive black holes to find that the "torus model" may be woefully inadequate when explaining what is actually going on. Most galaxies appear to contain a supermassiveblack hole in their cores. With masses in the realms of millions to billions of solar masses, these objects truly are the heavyweights of our Universe. With all this mass comes a powerful gravitational field that dominates galactic cores, pulling in any matter — stars, planets, dust, gas, possibly unlucky extraterrestrials — to the black hole's event horizon. [10 Strangest Black Hole Discoveries] Interactions between infalling matter and the supermassive black holes can generate huge quantities of energy, creating what are known as active galactic nuclei, making the effects of the black hole easy to observe. In the 1970s, astronomers developed a unified theory that could explain active supermassive black hole observations. The theory arose from the fact that some active black hole emissions could be easily seen by observatories while others seemed obscured by dust. To explain this, astronomers came up with the idea that supermassive black holes must be surrounded by a torus, or ring, of dusty material (as shown in the artistic rendering above). Therefore, given their random orientation as observed from Earth, some rings may appear "edge on" (thereby blocking our view of the black hole) or we may be observing the ring from above (revealing the black hole). Since this unified theory was suggested, it has generally matched observations of black holes and helped us understand how they influence the evolution of their host galaxies. However, new analyses of WISE data — a space telescope that surveyed the infrared sky twice for a little over a year until its primary mission was complete in February 2011 — has revealed a complication to the unified theory. As expected, after surveying 170,000 galaxies containing supermassive black holes at their cores, the WISE observations showed some black holes that could be seen, whereas others appeared obscured (in line with the torus model), but it also revealed a peculiar pattern. When looking at black holes inside massive galaxies that are clumped together as a part of galactic clusters, more supermassive black holes seemed to be obscured. This bias toward obscured black holes in large clusters cannot be accounted for if we just consider the unified theory. Why would supermassive black holes inside galaxies that are clumped in clusters be preferentially obscured by their dusty doughnut-shaped rings? "The main purpose of unification was to put a zoo of different kinds of active nuclei under a single umbrella," said post-doctorate astronomer and lead researcher Emilio Donoso, of the Instituto de Ciencias Astronómicas, de la Tierra y del Espacio in Argentina. "Now, that has become increasingly complex to do as we dig deeper into the WISE data." Donoso and his team's work indicates that some mechanism beyond the unified model is at work and they suggest dark matter may have a part to play. It is well known that "invisible" dark matter — a type of non-baryonic matter that pervades the entire universe — exerts a strong gravitational influence on galaxies and clusters of galaxies. It is also known that there is a vast, large-scale dark matter structure that forms a huge cosmic "web." At nodes in this dark matter web — known as halos — galaxies form and collect as clusters, apparently anchored in place by the gravitational oomph of dark matter halos. It is also known that some of the most massive supermassive black holes reside in the biggest, most massive clusters that, in turn, is the location of the biggest halos of dense dark matter. Could dark matter halos have a role in obscuring the clustered supermassive black holes from view? Is dark matter somehow adding more complexity to black hole torus? For now, we just don't know and further work is needed to understand this bias. But it is fascinating to think that the 'textbook' idea of an active black hole sporting a doughnut-shaped torus may need some reworking. This article was provided by Discovery News.
How to define a central angle and find the measure of its intercepted arc; how to describe the intercepted arcs of congruent chords. How to describe same side interior and same side exterior angles and their special properties. How to use a protractor to measure an angle. How to prove that opposite angles in a cyclic quadrilateral are congruent; how to prove that parallel lines create congruent arcs in a circle. How to identify a segment from the vertex angle in an isosceles triangle to the opposite side. How to relate corresponding altitudes, medians and angle bisectors in similar triangles. How to use the ASA and AAS shortcuts to determine the congruence of two triangles. How to define a tangent line; how to determine the angle a radius to a tangent forms. How to define isosceles triangles, their components and how to determine their properties. How to determine if two triangles in a circle are similar and how to prove that three similar triangles exist in a right triangle with an altitude. How to identify all of the special properties of parallelograms and use them to solve problems. How to construct the incenter using a compass and straightedge. How to identify a trapezoid and its special properties. How to define rotational symmetry and identify the degree of rotational symmetry of common regular polygons. How to identify a kite and its special properties. How to derive the area formula of a kite based on the rectangle formula; how to calculate the area of a rectangle using diagonal lengths. How to classify a triangle based on its side lengths and angle measures. How to find the measure of one angle in any equiangular or regular polygon. How to define the sine ratio and identify the sine of an angle in a right triangle. How to define the cosine ratio and identify the cosine of an angle in a right triangle.
If you find your head spinning when you think about all the addition strategies that you should be teaching, you are certainly not alone. Although teaching addition is one of the most important math concepts that we encounter, it comes with challenges. First of all – time. How do we find the time to do a really good job of teaching the different addition strategies so that our students possess excellent understanding, while also doing a really good job of teaching everything else in our overwhelming curriculum? Second, differentiation. All of our students learn at different speeds and in different ways. We can’t expect them all to learn the addition facts and strategies at the same time, but how do we ensure that each student is working to his full potential? One last big challenge is the balance between mental math strategies and memorization. We know that strategies are important. We want our students to UNDERSTAND number, rather that simply memorizing the facts. However, automaticity is important too! How can we reach this balance? Fortunately, there are specific strategies that we can teach to make addition easier for our students, and accessible for all of them. Before I begin – if you are looking for a resource where all of the work is done for you, you may be interested in The Addition Station, a self-paced, student-centered math station where students work through the basic addition facts and strategies, mastering each one as they go. Strategies are integrated in a strategic manner, ensuring that students build on their understanding progressively. See The Addition Station for Grades 1-2 HERE and The Addition Station for Grades 3-4 HERE. Alright, so let’s talk about addition strategies. Strategies are ESSENTIAL, for all operations. We want our students to be able to think flexibly about numbers, and use strategies naturally. This means that understanding is key. Automaticity (quick recall) will naturally follow. Here are some effective strategies for addition: Plus 1, 2, 3 and Extensions In younger grades, we begin with the Plus 1, 2, and 3 facts. We can teach Plus 1 as 1 more, Plus 2 as 2 more, and Plus 3 as 3 more. As our students are ready for more of a challenge, we can extend these facts into the tens, hundreds, and even thousands. For example, the fact 7+1 can be extended to 70+10, 700+100, or 7000+1000. Teach your students to look for familiar facts in these bigger problems, so that when they need to solve a fact like 50+20, they think, “I know that 5+2=7, so 50+20=70.” During the extensions, be sure to emphasize place value. For example, we can think of 500+200 as 5 groups of 100 plus 2 groups of 100 to make 7 groups of 100. If students have been working with Plus 1, 2, and 3, they have technically already been working with the counting on strategy. Counting On is an introductory addition strategy that should only be used to add 1, 2, 3, or 4 to a number. Beyond this it gets confusing and can cause errors. To count on, we begin with the higher number and count on. For example, for 17+3, we think, “17…18, 19, 20.” For 2+34, we start with 34 and count on: “34…35, 36.” Dot patterns, ten frames, and number lines are all excellent tools for counting on. Read more about counting on, and download some free printables to help you HERE. Extending the Doubles and Near Doubles The doubles are typically facts that become automatic early on. The near doubles are facts like 4+5, where we encourage students to think, “I know that 4+4 is 8, and 1 more is 9.” We can also extend these facts. For example, when a student is faced with 30+30, he can think, “I know that 3+3 is 6, so 30+30 is 60.” Again in this level, we encourage students to think in groups of 10, 100, or 1000. For example, 200+200 can be thought of as 2 groups of 100 plus 2 groups of 100. Read more about extending the doubles and near doubles, and download some free printables to help you HERE. Plus 7, 8, and 9 When we add 7, 8, or 9 to a number there are a couple of different approaches that we can use. First of all, we can add 10 and then take some away. Alternatively, we can make a 10 and then add the rest. I’ve discussed both of these approaches in detail in THIS BLOG POST. You’ll also find some free printables to help you! Left-to-right addition (also known as front-end addition or the partial sums method) is one of the most powerful mental math strategies for teaching addition of 2 or 3-digit numbers. However, many people are confused by why it is important and why it can be more effective than traditional vertical addition. With left-to-right addition we add from left to right. So in a two-digit equation we add the tens first and then the ones. For example for 25+34 we first add 20+30 to make 50, then 5+4 to make 9, and then 50+9 to make the final sum of 59. For a detailed explanation and rationalization of this strategy, as well as free printables to help you teach it, please see THIS POST. Using Friendly Numbers A friendly number is a number that is easy to work with. For example, multiples of 10 are “friendly” because they are easy to work with when we add or subtract. When we use the “friendly number” strategy for addition, it helps us work with big numbers. This is because we are essentially breaking the equation up into more manageable parts. We begin by getting to a friendly number, which is typically a multiple of 10, 100, or 100 – depending on the numbers that we are working with. Then we add on the remainder. For example, for the equation 27+9, we could first get to the friendly number 30 by adding 3, and then add the remaining 6 to make 36. For a complete, detailed explanation of the friendly number strategy for addition, as well as some free printables to help you, please see THIS POST. Breaking Up An Addend Breaking up, or decomposing, an addend is a fantastic mental math strategy for addition that can be used in many different circumstances. This strategy involves breaking up one of the numbers in an equation into more manageable parts. Like many other mental math strategies, this encourages students to think flexibly and to manipulate numbers in different ways. This is the big goal of mental math! For a detailed explanation of this strategy as well as free printables to help you teach it, please see THIS POST. Compensation is a mental math strategy for multi-digit addition that involves adjusting one of the addends to make the equation easier to solve. Some students may prefer this strategy as an alternative to left-to-right addition or the breaking up the second number strategy. Compensation is a useful strategy for making equations easier to solve. More importantly, it encourages students to think flexibly about numbers. Let’s solve the equation 34+49 using the compensation strategy. First, since 49 is so close to 50, we will add 34+50. This is easier to solve. Then, since we added one extra to the original equation, we have to subtract one from the final answer. To see a more detailed explanation of this strategy, and download some free printables to teach it, please see THIS POST. Ready to get really strategic with your approach to teaching math facts? - Implement The Addition Station – a self-paced, student-centered program for the basic addition facts and strategies (see Grades 1-2 HERE and Grades 3-4 HERE). - Try out Addition Strategies Task Cards as a way for your students to practice each addition strategy in isolation. See the full bundle for 1st Grade, 2nd grade, 3rd grade, or 4th grade.
Measurement in Standard Units Young scholars choose the proper measurement unit and the corresponding measurement tool needed for a particular measurement in length, volume, and weight/mass. 6 Views 35 Downloads - Activities & Projects - Graphics & Images - Handouts & References - Lab Resources - Learning Games - Lesson Plans - Primary Sources - Printables & Templates - Professional Documents - Study Guides - Graphic Organizers - Writing Prompts - Constructed Response Items - AP Test Preps - Lesson Planet Articles - Interactive Whiteboards - All Resource Types - Show All See similar resources: Measuring Mass Using Non-Standard Units of Measurement Students explore mass measurements. In this non-standard unit measurement lesson, students use non-standard units of measurement to find the weight of different objects. They work in small groups and complete a variety of measuring... 1st - 2nd Math How Tall is the Gingerbread Man? The gingerbread man has finally been caught, now find out how tall he is with this measurement activity. Using Unifix® cubes, marshmallows, buttons, and Cheerios, young mathematicians measure and record the height of the gingerbread man... Pre-K - 2nd Math CCSS: Adaptable Convert Between Units: Using a T-Chart Learning to convert between different units of measurement can be a challenge for many young mathematicians. The final video in this series supports learners with this topic by clearly modeling how T-charts and repeated addition are used... 5 mins 2nd - 6th Math CCSS: Designed Cover the Area of a Shape Using Square Units This is the first of a series of four lessons designed to examine and understand the concept of area. A review of area starts the lesson, along with a discussion of different attributes that can describe a rectangle. The term square unit... 3 mins 2nd - 4th Math CCSS: Designed
\|Part of a series on\| Readability is the ease with which a reader can understand a written text. In natural language, the readability of text depends on its content (the complexity of its vocabulary and syntax) and its presentation (such as typographic aspects that affect legibility, like font size, line height, character spacing, and line length). Researchers have used various factors to measure readability, such as: - Speed of perception - Perceptibility at a distance - Perceptibility in peripheral vision - Reflex blink technique - Rate of work (reading speed) - Eye movements - Fatigue in reading - Cognitively-motivated features - Word difficulty - N-gram analysis - Semantic Richness Higher readability eases reading effort and speed for any reader, but it makes a larger difference for those who do not have high reading comprehension. Readability exists in both natural language and programming languages though in different forms. In programming, things such as programmer comments, choice of loop structure, and choice of names can determine the ease with which humans can read computer program code. Numeric readability metrics (also known as readability tests or readability formulas) for natural language tend to use simple measures like word length (by letter or syllable), sentence length, and sometimes some measure of word frequency. They can be built into word processors, can score documents, paragraphs, or sentences, and are a much cheaper and faster alternative to a readability survey involving human readers. They are faster to calculate than more accurate measures of syntactic and semantic complexity. In some cases they are used to estimate appropriate grade level. People have defined readability in various ways, e.g., in: The Literacy Dictionary, Jeanne Chall and Edgar Dale, G. Harry McLaughlin, William DuBay.[further explanation needed] Easy reading helps learning and enjoyment, and can save money. Much research has focused on matching prose to reading skill, resulting in formulas for use in research, government, teaching, publishing, the military, medicine, and business. Readability and newspaper readership Several studies in the 1940s showed that even small increases in readability greatly increases readership in large-circulation newspapers. In 1947, Donald Murphy of Wallace's Farmer used a split-run edition to study the effects of making text easier to read. He found that reducing from the 9th to the 6th-grade reading level increased readership by 43% for an article on 'nylon'. The result was a gain of 42,000 readers in a circulation of 275,000. He also found a 60% increase in readership for an article on corn, with better responses from people under 35. Wilber Schramm interviewed 1,050 newspaper readers. He found that an easier reading style helps to determine how much of an article is read. This was called reading persistence, depth, or perseverance. He also found that people will read less of long articles than of short ones. A story nine paragraphs long will lose 3 out of 10 readers by the fifth paragraph. A shorter story will lose only two. Schramm also found that the use of subheads, bold-face paragraphs, and stars to break up a story actually lose readers. A study in 1947 by Melvin Lostutter showed that newspapers generally were written at a level five years above the ability of average American adult readers. The reading ease of newspaper articles was not found to have much connection with the education, experience, or personal interest of the journalists writing the stories. It instead had more to do with the convention and culture of the industry. Lostutter argued for more readability testing in newspaper writing. Improved readability must be a "conscious process somewhat independent of the education and experience of the staffs writers." A study by Charles Swanson in 1948 showed that better readability increases the total number of paragraphs read by 93% and the number of readers reading every paragraph by 82%. In 1948, Bernard Feld did a study of every item and ad in the Birmingham News of 20 November 1947. He divided the items into those above the 8th-grade level and those at the 8th grade or below. He chose the 8th-grade breakpoint, as that was determined to be the average reading level of adult readers. An 8th-grade text "...will reach about 50% of all American grown-ups," he wrote. Among the wire-service stories, the lower group got two-thirds more readers, and among local stories, 75% more readers. Feld also believed in drilling writers in Flesch's clear-writing principles. Both Rudolf Flesch and Robert Gunning worked extensively with newspapers and the wire services in improving readability. Mainly through their efforts in a few years, the readability of US newspapers went from the 16th to the 11th-grade level, where it remains today. The two publications with the largest circulations, TV Guide (13 million) and Reader's Digest (12 million), are written at the 9th-grade level. The most popular novels are written at the 7th-grade level. This supports the fact that the average adult reads at the 9th-grade level. It also shows that, for recreation, people read texts that are two grades below their actual reading level. The George Klare studies George Klare and his colleagues looked at the effects of greater reading ease on Air Force recruits. They found that more readable texts resulted in greater and more complete learning. They also increased the amount read in a given time, and made for easier acceptance. Other studies by Klare showed how the reader's skills, prior knowledge, interest, and motivation affect reading ease. In the 1880s, English professor L. A. Sherman found that the English sentence was getting shorter. In Elizabethan times, the average sentence was 50 words long. In his own time, it was 23 words long. Sherman's work established that: - Literature is a subject for statistical analysis. - Shorter sentences and concrete terms help people to make sense of what is written. - Speech is easier to understand than text. - Over time, text becomes easier if it is more like speech. Sherman wrote: "Literary English, in short, will follow the forms of standard spoken English from which it comes. No man should talk worse than he writes, no man should write better than he should talk.... The oral sentence is clearest because it is the product of millions of daily efforts to be clear and strong. It represents the work of the race for thousands of years in perfecting an effective instrument of communication." In 1889 in Russia, the writer Nikolai A. Rubakin published a study of over 10,000 texts written by everyday people. From these texts, he took 1,500 words he thought most people understood. He found that the main blocks to comprehension are unfamiliar words and long sentences. Starting with his own journal at the age of 13, Rubakin published many articles and books on science and many subjects for the great numbers of new readers throughout Russia. In Rubakin's view, the people were not fools. They were simply poor and in need of cheap books, written at a level they could grasp. In 1921, Harry D. Kitson published The Mind of the Buyer, one of the first books to apply psychology to marketing. Kitson's work showed that each type of reader bought and read their own type of text. On reading two newspapers and two magazines, he found that short sentence length and short word length were the best contributors to reading ease. The earliest reading ease assessment is the subjective judgment termed text leveling. Formulas do not fully address the various content, purpose, design, visual input, and organization of a text. Text leveling is commonly used to rank the reading ease of texts in areas where reading difficulties are easy to identify, such as books for young children. At higher levels, ranking reading ease becomes more difficult, as individual difficulties become harder to identify. This has led to better ways to assess reading ease. Vocabulary frequency lists In the 1920s, the scientific movement in education looked for tests to measure students' achievement to aid in curriculum development. Teachers and educators had long known that, to improve reading skill, readers—especially beginning readers—need reading material that closely matches their ability. University-based psychologists did much of the early research, which was later taken up by textbook publishers. Educational psychologist Edward Thorndike of Columbia University noted that, in Russia and Germany, teachers used word frequency counts to match books to students. Word skill was the best sign of intellectual development, and the strongest predictor of reading ease. In 1921, Thorndike published Teachers Word Book, which contained the frequencies of 10,000 words. It made it easier for teachers to choose books that matched class reading skills. It also provided a basis for future research on reading ease. Until computers came along, word frequency lists were the best aids for grading reading ease of texts. In 1981 the World Book Encyclopedia listed the grade levels of 44,000 words. Early children's readability formulas In 1923, Bertha A. Lively and Sidney L. Pressey published the first reading ease formula. They were concerned that junior high school science textbooks had so many technical words. They felt that teachers spent all class time explaining these words. They argued that their formula would help to measure and reduce the "vocabulary burden" of textbooks. Their formula used five variable inputs and six constants. For each thousand words, it counted the number of unique words, the number of words not on the Thorndike list, and the median index number of the words found on the list. Manually, it took three hours to apply the formula to a book. After the Lively–Pressey study, people looked for formulas that were more accurate and easier to apply. By 1980, over 200 formulas were published in different languages. In 1928, Carleton Washburne and Mabel Vogel created the first modern readability formula. They validated it by using an outside criterion, and correlated .845 with test scores of students who read and liked the criterion books. It was also the first to introduce the variable of interest to the concept of readability. Between 1929 and 1939, Alfred Lewerenz of the Los Angeles School District published several new formulas. In 1934, Edward Thorndike published his formula. He wrote that word skills can be increased if the teacher introduces new words and repeats them often. In 1939, W.W. Patty and W. I Painter published a formula for measuring the vocabulary burden of textbooks. This was the last of the early formulas that used the Thorndike vocabulary-frequency list. Early adult readability formulas During the recession of the 1930s, the U.S. government invested in adult education. In 1931, Douglas Waples and Ralph Tyler published What Adults Want to Read About. It was a two-year study of adult reading interests. Their book showed not only what people read but what they would like to read. They found that many readers lacked suitable reading materials: they would have liked to learn but the reading materials were too hard for them. Lyman Bryson of Teachers College, Columbia University found that many adults had poor reading ability due to poor education. Even though colleges had long tried to teach how to write in a clear and readable style, Bryson found that it was rare. He wrote that such language is the result of a "...discipline and artistry that few people who have ideas will take the trouble to achieve... If simple language were easy, many of our problems would have been solved long ago." Bryson helped set up the Readability Laboratory at the college. Two of his students were Irving Lorge and Rudolf Flesch. In 1934, Ralph Ojemann investigated adult reading skills, factors that most directly affect reading ease, and causes of each level of difficulty. He did not invent a formula, but a method for assessing the difficulty of materials for parent education. He was the first to assess the validity of this method by using 16 magazine passages tested on actual readers. He evaluated 14 measurable and three reported factors that affect reading ease. Ojemann emphasized the reported features, such as whether the text was coherent or unduly abstract. He used his 16 passages to compare and judge the reading ease of other texts, a method now called scaling. He showed that even though these factors cannot be measured, they cannot be ignored. Also in 1934, Ralph Tyler and Edgar Dale published the first adult reading ease formula based on passages on health topics from a variety of textbooks and magazines. Of 29 factors that are significant for young readers, they found ten that are significant for adults. They used three of these in their formula. In 1935, William S. Gray of the University of Chicago and Bernice Leary of Xavier College in Chicago published What Makes a Book Readable, one of the most important books in readability research. Like Dale and Tyler, they focused on what makes books readable for adults of limited reading ability. Their book included the first scientific study of the reading skills of American adults. The sample included 1,690 adults from a variety of settings and regions. The test used a number of passages from newspapers, magazines, and books—as well as a standard reading test. They found a mean grade score of 7.81 (eighth month of the seventh grade). About one-third read at the 2nd to 6th-grade level, one-third at the 7th to 12th-grade level, and one-third at the 13th–17th grade level. The authors emphasized that one-half of the adult population at that time lacked suitable reading materials. They wrote, "For them, the enriching values of reading are denied unless materials reflecting adult interests are adapted to their needs." The poorest readers, one-sixth of the adult population, need "simpler materials for use in promoting functioning literacy and in establishing fundamental reading habits." Gray and Leary then analyzed 228 variables that affect reading ease and divided them into four types: They found that content was most important, followed closely by style. Third was format, followed closely by organization. They found no way to measure content, format, or organization—but they could measure variables of style. Among the 17 significant measurable style variables, they selected five to create a formula: - Average sentence length - Number of different hard words - Number of personal pronouns - Percentage of unique words - Number of prepositional phrases Their formula had a correlation of .645 with comprehension as measured by reading tests given to about 800 adults. In 1939, Irving Lorge published an article that reported other combinations of variables that indicate difficulty more accurately than the ones Gray and Leary used. His research also showed that, "The vocabulary load is the most important concomitant of difficulty." In 1944, Lorge published his Lorge Index, a readability formula that used three variables and set the stage for simpler and more reliable formulas that followed. By 1940, investigators had: - Successfully used statistical methods to analyze reading ease - Found that unusual words and sentence length were among the first causes of reading difficulty - Used vocabulary and sentence length in formulas to predict reading ease Popular readability formulas The Flesch formulas In 1943, Rudolf Flesch published his PhD dissertation, Marks of a Readable Style, which included a readability formula to predict the difficulty of adult reading material. Investigators in many fields began using it to improve communications. One of the variables it used was personal references, such as names and personal pronouns. Another variable was affixes. In 1948, Flesch published his Reading Ease formula in two parts. Rather than using grade levels, it used a scale from 0 to 100, with 0 equivalent to the 12th grade and 100 equivalent to the 4th grade. It dropped the use of affixes. The second part of the formula predicts human interest by using personal references and the number of personal sentences. The new formula correlated 0.70 with the McCall-Crabbs reading tests. The original formula is: - Reading Ease score = 206.835 − (1.015 × ASL) − (84.6 × ASW) - Where: ASL = average sentence length (number of words divided by number of sentences) - ASW = average word length in syllables (number of syllables divided by number of words) Publishers discovered that the Flesch formulas could increase readership up to 60%. Flesch's work also made an enormous impact on journalism. The Flesch Reading Ease formula became one of the most widely used, tested, and reliable readability metrics. In 1951, Farr, Jenkins, and Patterson simplified the formula further by changing the syllable count. The modified formula is: - New reading ease score = 1.599nosw − 1.015sl − 31.517 - Where: nosw = number of one-syllable words per 100 words and - sl = average sentence length in words. In 1975, in a project sponsored by the U.S. Navy, the Reading Ease formula was recalculated to give a grade-level score. The new formula is now called the Flesch–Kincaid grade-level formula. The Flesch–Kincaid formula is one of the most popular and heavily tested formulas. It correlates 0.91 with comprehension as measured by reading tests. The Dale–Chall formula Edgar Dale, a professor of education at Ohio State University, was one of the first critics of Thorndike's vocabulary-frequency lists. He claimed that they did not distinguish between the different meanings that many words have. He created two new lists of his own. One, his "short list" of 769 easy words, was used by Irving Lorge in his formula. The other was his "long list" of 3,000 easy words, which were understood by 80% of fourth-grade students. However, one has to extend the word lists by regular plurals of nouns, regular forms of the past tense of verbs, progressive forms of verbs etc. In 1948, he incorporated this list into a formula he developed with Jeanne S. Chall, who later founded the Harvard Reading Laboratory. To apply the formula: - Select several 100-word samples throughout the text. - Compute the average sentence length in words (divide the number of words by the number of sentences). - Compute the percentage of words NOT on the Dale–Chall word list of 3,000 easy words. - Compute this equation from 1948: - Raw score = 0.1579(PDW) + 0.0496(ASL) if the percentage of PDW is less than 5%, otherwise compute - Raw score = 0.1579(PDW) + 0.0496(ASL) + 3.6365 - Raw score = uncorrected reading grade of a student who can answer one-half of the test questions on a passage. - PDW = Percentage of difficult words not on the Dale–Chall word list. - ASL = Average sentence length Finally, to compensate for the "grade-equivalent curve", apply the following chart for the Final Score: \|Raw score\|\|Final score\| \|4.9 and below\|\|Grade 4 and below\| \|9.0–9.9\|\|Grades 13–15 (college)\| \|10 and above\|\|Grades 16 and above.\| Correlating 0.93 with comprehension as measured by reading tests, the Dale–Chall formula is the most reliable formula and is widely used in scientific research. In 1995, Dale and Chall published a new version of their formula with an upgraded word list, the New Dale–Chall readability formula. Its formula is: Raw score = 64 - 0.95 (PDW) - 0.69 (ASL) The Gunning fog formula In the 1940s, Robert Gunning helped bring readability research into the workplace. In 1944, he founded the first readability consulting firm dedicated to reducing the "fog" in newspapers and business writing. In 1952, he published The Technique of Clear Writing with his own Fog Index, a formula that correlates 0.91 with comprehension as measured by reading tests. The formula is one of the most reliable and simplest to apply: - Grade level= 0.4 * ( (average sentence length) + (percentage of Hard Words) ) - Where: Hard Words = words with more than two syllables. Fry readability graph In 1963, while teaching English teachers in Uganda, Edward Fry developed his Readability Graph. It became one of the most popular formulas and easiest to apply. The Fry Graph correlates 0.86 with comprehension as measured by reading tests. McLaughlin's SMOG formula Harry McLaughlin determined that word length and sentence length should be multiplied rather than added as in other formulas. In 1969, he published his SMOG (Simple Measure of Gobbledygook) formula: - SMOG grading = 3 + √polysyllable count. - Where: polysyllable count = number of words of more than two syllables in a sample of 30 sentences. The SMOG formula correlates 0.88 with comprehension as measured by reading tests. It is often recommended for use in healthcare. The FORCAST formula In 1973, a study commissioned by the US military of the reading skills required for different military jobs produced the FORCAST formula. Unlike most other formulas, it uses only a vocabulary element, making it useful for texts without complete sentences. The formula satisfied requirements that it would be: - Based on Army-job reading materials. - Suitable for the young adult-male recruits. - Easy enough for Army clerical personnel to use without special training or equipment. The formula is: - Grade level = 20 − (N / 10) - Where N = number of single-syllable words in a 150-word sample. The FORCAST formula correlates 0.66 with comprehension as measured by reading tests. The Golub Syntactic Density Score The Golub Syntactic Density Score was developed by Lester Golub in 1974. It is among a smaller subset of readability formulas that concentrate on the syntactic features of a text. To calculate the reading level of a text, a sample of several hundred words is taken from the text. The number of words in the sample is counted, as are the number of T-units. A T-unit is defined as an independent clause and any dependent clauses attached to it. Other syntactical units are then counted and entered into the following table: 1. Words/T-unit .95 X _________ ___ 2. Subordinate clauses/T-unit .90 X _________ ___ 3. Main clause word length (mean) .20 X _________ ___ 4. Subordinate clause length (mean) .50 X _________ ___ 5. Number of Modals (will, shall, can, may, must, would...) .65 X _________ ___ 6. Number of Be and Have forms in the auxiliary .40 X _________ ___ 7. Number of Prepositional Phrases .75 X _________ ___ 8. Number of Possessive nouns and pronouns .70 X _________ ___ 9. Number of Adverbs of Time (when, then, once, while...) .60 X _________ ___ 10. Number of gerunds, participles, and absolutes Phrases .85 X _________ ___ Users add the numbers in the right hand column and divide the total by the number of T-units. Finally, the quotient is entered into the following table to arrive at a final readability score. Measuring coherence and organization For centuries, teachers and educators have seen the importance of organization, coherence, and emphasis in good writing. Beginning in the 1970s, cognitive theorists began teaching that reading is really an act of thinking and organization. The reader constructs meaning by mixing new knowledge into existing knowledge. Because of the limits of the reading ease formulas, some research looked at ways to measure the content, organization, and coherence of text. Although this did not improve the reliability of the formulas, their efforts showed the importance of these variables in reading ease. Studies by Walter Kintch and others showed the central role of coherence in reading ease, mainly for people learning to read. In 1983, Susan Kemper devised a formula based on physical states and mental states. However, she found this was no better than word familiarity and sentence length in showing reading ease. Bonnie Meyer and others tried to use organization as a measure of reading ease. While this did not result in a formula, they showed that people read faster and retain more when the text is organized in topics. She found that a visible plan for presenting content greatly helps readers to assess a text. A hierarchical plan shows how the parts of the text are related. It also aids the reader in blending new information into existing knowledge structures. Bonnie Armbruster found that the most important feature for learning and comprehension is textual coherence, which comes in two types: - Global coherence, which integrates high-level ideas as themes in an entire section, chapter, or book. - Local coherence, which joins ideas within and between sentences. Armbruster confirmed Kintsch's finding that coherence and structure are more help for younger readers. R. C. Calfee and R. Curley built on Bonnie Meyer's work and found that an unfamiliar underlying structure can make even simple text hard to read. They brought in a graded system to help students progress from simpler story lines to more advanced and abstract ones. Many other studies looked at the effects on reading ease of other text variables, including: - Image words, abstraction, direct and indirect statements, types of narration and sentences, phrases, and clauses; - Difficult concepts; - Idea density; - Human interest; - Active and passive voice; - Structural cues; - The use of images; - Diagrams and line graphs; - Fonts and layout; - Document age. Advanced readability formulas The John Bormuth formulas John Bormuth of the University of Chicago looked at reading ease using the new Cloze deletion test developed by Wilson Taylor. His work supported earlier research including the degree of reading ease for each kind of reading. The best level for classroom "assisted reading" is a slightly difficult text that causes a "set to learn", and for which readers can correctly answer 50% of the questions of a multiple-choice test. The best level for unassisted reading is one for which readers can correctly answer 80% of the questions. These cutoff scores were later confirmed by Vygotsky and Chall and Conard. Among other things, Bormuth confirmed that vocabulary and sentence length are the best indicators of reading ease. He showed that the measures of reading ease worked as well for adults as for children. The same things that children find hard are the same for adults of the same reading levels. He also developed several new measures of cutoff scores. One of the most well known was the Mean Cloze Formula, which was used in 1981 to produce the Degree of Reading Power system used by the College Entrance Examination Board. The Lexile framework In 1988, Jack Stenner and his associates at MetaMetrics, Inc. published a new system, the Lexile Framework, for assessing readability and matching students with appropriate texts. The Lexile framework uses average sentence length, and average word frequency in the American Heritage Intermediate Corpus to predict a score on a 0–2000 scale. The AHI Corpus includes five million words from 1,045 published works often read by students in grades three to nine. The Lexile Book Database has more than 100,000 titles from more than 450 publishers. By knowing a student's Lexile score, a teacher can find books that match his or her reading level. ATOS readability formula for books In 2000, researchers of the School Renaissance Institute and Touchstone Applied Science Associates published their Advantage-TASA Open Standard (ATOS) Reading ease Formula for Books. They worked on a formula that was easy to use and that could be used with any texts. The project was one of the widest reading ease projects ever. The developers of the formula used 650 normed reading texts, 474 million words from all the text in 28,000 books read by students. The project also used the reading records of more than 30,000 who read and were tested on 950,000 books. They found that three variables give the most reliable measure of text reading ease: - words per sentence - average grade level of words - characters per word They also found that: - To help learning, the teacher should match book reading ease with reading skill. - Reading often helps with reading gains. - For reading alone below the 4th grade, the best learning gain requires at least 85% comprehension. - Advanced readers need 92% comprehension for independent reading. - Book length can be a good measure of reading ease. - Feedback and interaction with the teacher are the most important factors in reading. CohMetrix psycholinguistics measurements Coh-Metrix can be used in many different ways to investigate the cohesion of the explicit text and the coherence of the mental representation of the text. "Our definition of cohesion consists of characteristics of the explicit text that play some role in helping the reader mentally connect ideas in the text." The definition of coherence is the subject of much debate. Theoretically, the coherence of a text is defined by the interaction between linguistic representations and knowledge representations. While coherence can be defined as characteristics of the text (i.e., aspects of cohesion) that are likely to contribute to the coherence of the mental representation, Coh-Metrix measurements provide indices of these cohesion characteristics. - Automated readability index (1967) - Linsear Write Raygor readability estimate (1977) - Spache readability formula (1952) Artificial Intelligence (AI) approach Unlike the traditional readability formulas, artificial intelligence approaches to readability assessment (also known as Automatic Readability Assessment) incorporate myriad linguistic features and construct statistical prediction models to predict text readability. These approaches typically consist of three steps: 1. a training corpus of individual texts, 2. a set of linguistic features to be computed from each text, and 3. a machine learning model to predict the readability, using the computed linguistic feature values. In 2012, Sowmya Vajjala at the University of Tübingen created the WeeBit corpus by combining educational articles from the Weekly Reader website and BBC Bitesize website, which provide texts for different age groups. In total, there are 3125 articles that are divided into five readability levels (from age 7 to 16). Weebit corpus has been used in several AI-based readability assessment research. Wei Xu (University of Pennsylvania), Chris Callison-Burch (University of Pennsylvania), and Courtney Napoles (Johns Hopkins University) introduced the Newsela corpus to the academic field in 2015. The corpus is a collection of thousands of news articles professionally leveled to different reading complexities by professional editors at Newsela. The corpus was originally introduced for text simplification research, but was also used for text readability assessment. Semantic or Advanced Semantic Advanced semantic or semantic features' influence on text readability was pioneered by Bruce W. Lee during his study at the (University of Pennsylvania), in 2021. Whilst introducing his features hybridization method, he also explored handcrafted advanced semantic features which aim to measure the amount of knowledge contained in a given text. - Semantic Richness : - Semantic Clarity : - Semantic Noise : where the count of discovered topics (n) and topic probability (p) The type-token ratio is one of the features that are often used to captures the lexical richness, which is a measure of vocabulary range and diversity. To measure the lexical difficulty of a word, the relative frequency of the word in a representative corpus like the Corpus of Contemporary American English (COCA) is often used. Below includes some examples for lexico-semantic features in readability assessment. - Average number of syllables per word - Out-of-vocabulary rate, in comparison to the full corpus - Type-token ratio: the ratio of unique terms to total terms observed - Ratio of function words, in comparison to the full corpus - Ratio of pronouns, in comparison to the full corpus - Language model perplexity (comparing the text to generic or genre-specific models) In addition, Lijun Feng pioneered the cognitively-motivated features (mostly lexical) in 2009. This was during her doctorate study at the City University of New York (CUNY). The cognitively-motivated features were originally designed for adults with intellectual disability, but was proved to improve readability assessment accuracy in general. Cognitively-motivated features, in combination with a logistic regression model, can correct the average error of Flesch–Kincaid grade-level by more than 70%. The newly discovered features by Feng include: - Number of lexical chains in document - Average number of unique entities per sentence - Average number of entity mentions per sentence - Total number of unique entities in document - Total number of entity mentions in document - Average lexical chain length - Average lexical chain span Syntactic complexity is correlated with longer processing times in text comprehension. It is common to use a rich set of these syntactic features to predict the readability of a text. The more advanced variants of syntactic readability features are frequently computed from parse tree. Emily Pitler (University of Pennsylvania) and Ani Nenkova (University of Pennsylvania) are considered pioneers in evaluating the parse-tree syntactic features and making it widely used in readability assessment. Some examples include: - Average sentence length - Average parse tree height - Average number of noun phrases per sentence - Average number of verb phrases per sentence Using the readability formulas The accuracy of readability formulas increases when finding the average readability of a large number of works. The tests generate a score based on characteristics such as statistical average word length (which is used as an unreliable proxy for semantic difficulty; sometimes word frequency is taken into account) and sentence length (as an unreliable proxy for syntactic complexity) of the work. Most experts agree that simple readability formulas like Flesch–Kincaid grade-level can be highly misleading. Even though the traditional features like the average sentence length have high correlation with reading difficulty, the measure of readability is much more complex. The artificial intelligence, data-driven approach (see above) was studied to tackle this shortcoming. Writing experts have warned that an attempt to simplify the text only by changing the length of the words and sentences may result in text that is more difficult to read. All the variables are tightly related. If one is changed, the others must also be adjusted, including approach, voice, person, tone, typography, design, and organization. Writing for a class of readers other than one's own is very difficult. It takes training, method, and practice. Among those who are good at this are writers of novels and children's books. The writing experts all advise that, besides using a formula, observe all the norms of good writing, which are essential for writing readable texts. Writers should study the texts used by their audience and their reading habits. This means that for a 5th-grade audience, the writer should study and learn good quality 5th-grade materials. - Asemic writing - Plain language - Accessible publishing - George R. Klare - William S. 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West Nile fever \|West Nile fever\| \|West Nile virus\| \|Symptoms\|\|None, fever, headache, vomiting, rash\| \|Usual onset\|\|2 to 14 days after exposure\| \|Duration\|\|Weeks to months\| \|Causes\|\|West Nile virus spread by mosquito\| \|Diagnostic method\|\|Based on symptoms and blood tests\| \|Prevention\|\|Reducing mosquitoes, preventing mosquito bites\| \|Treatment\|\|Supportive care (pain medication)\| \|Prognosis\|\|10% risk of death among those seriously affected\| West Nile fever is an infection by the West Nile virus, which is typically spread by mosquitoes. In about 80% of infections people have few or no symptoms. About 20% of people develop a fever, headache, vomiting, or a rash. In less than 1% of people, encephalitis or meningitis occurs, with associated neck stiffness, confusion, or seizures. Recovery may take weeks to months; the risk of death among those in whom the nervous system is affected is about 10%. West Nile virus (WNV) is usually spread by infected mosquitoes. Mosquitoes become infected when they feed on infected birds, which often carry the disease. Rarely the virus is spread through blood transfusions, organ transplants, or from mother to baby during pregnancy, delivery, or breastfeeding, it otherwise does not spread directly between people. Risks for severe disease include being over 60 years old and having other health problems. Diagnosis is typically based on symptoms and blood tests. There is no human vaccine; the best way to reduce the risk of infection is to avoid mosquito bites. Mosquito populations may be reduced by eliminating standing pools of water, such as in old tires, buckets, gutters, and swimming pools; when mosquitoes cannot be avoided, mosquito repellent, window screens, and mosquito nets reduce the likelihood of being bitten. There is no specific treatment for the disease; pain medications may reduce symptoms. The virus was discovered in Uganda in 1937, and was first detected in North America in 1999. WNV has occurred in Europe, Africa, Asia, Australia, and North America. In the United States thousands of cases are reported a year, with most occurring in August and September, it can occur in outbreaks of disease. Severe disease may also occur in horses, for which a vaccine is available. A surveillance system in birds is useful for early detection of a potential human outbreak. Signs and symptoms About 80% of those infected with West Nile virus (WNV) show no symptoms and go unreported. About 20% of infected people develop symptoms; these vary in severity, and begin 3 to 14 days after being bitten. Most people with mild symptoms of WNV recover completely, though fatigue and weakness may last for weeks or months. Symptoms may range from mild, such as fever, to severe, such as paralysis and meningitis. A severe infection can last weeks and can, rarely, cause permanent brain damage. Death may ensue if the central nervous system is affected. Medical conditions such as cancer and diabetes, and age over 60 years, increase the risk of developing severe symptoms. Headache can be a prominent symptom of WNV fever, meningitis, encephalitis, meningoencephalitis, and it may or may not be present in poliomyelitis-like syndrome. Thus, headache is not a useful indicator of neuroinvasive disease. - West Nile fever (WNF), which occurs in 20 percent of cases, is a febrile syndrome that causes flu-like symptoms. Most characterizations of WNF describe it as a mild, acute syndrome lasting 3 to 6 days after symptom onset. Systematic follow-up studies of patients with WNF have not been done, so this information is largely anecdotal. Possible symptoms include high fever, headache, chills, excessive sweating, weakness, fatigue, swollen lymph nodes, drowsiness, pain in the joints and flu-like symptoms. There may be gastrointestinal symptoms including nausea, vomiting, loss of appetite, and diarrhea. Fewer than one-third of patients develop a rash. - West Nile neuroinvasive disease (WNND), which occurs in less than 1 percent of cases, is when the virus infects the central nervous system resulting in meningitis, encephalitis, meningoencephalitis or a poliomyelitis-like syndrome. Many patients with WNND have normal neuroimaging studies, although abnormalities may be present in various cerebral areas including the basal ganglia, thalamus, cerebellum, and brainstem. - West Nile virus encephalitis (WNE) is the most common neuroinvasive manifestation of WNND. WNE presents with similar symptoms to other viral encephalitis with fever, headaches, and altered mental status. A prominent finding in WNE is muscular weakness (30 to 50 percent of patients with encephalitis), often with lower motor neuron symptoms, flaccid paralysis, and hyporeflexia with no sensory abnormalities. - West Nile meningitis (WNM) usually involves fever, headache, stiff neck and pleocytosis, an increase of white blood cells in cerebrospinal fluid. Changes in consciousness are not usually seen and are mild when present. - West Nile meningoencephalitis is inflammation of both the brain (encephalitis) and meninges (meningitis). - West Nile poliomyelitis (WNP), an acute flaccid paralysis syndrome associated with WNV infection, is less common than WNM or WNE. This syndrome is generally characterized by the acute onset of asymmetric limb weakness or paralysis in the absence of sensory loss. Pain sometimes precedes the paralysis; the paralysis can occur in the absence of fever, headache, or other common symptoms associated with WNV infection. Involvement of respiratory muscles, leading to acute respiratory failure, sometimes occurs. - West-Nile reversible paralysis, Like WNP, the weakness or paralysis is asymmetric. Reported cases have been noted to have an initial preservation of deep tendon reflexes, which is not expected for a pure anterior horn involvement. Disconnect of upper motor neuron influences on the anterior horn cells possibly by myelitis or glutamate excitotoxicity have been suggested as mechanisms; the prognosis for recovery is excellent. - Nonneurologic complications of WNV infection that may rarely occur include fulminant hepatitis, pancreatitis, myocarditis, rhabdomyolysis, orchitis, nephritis, optic neuritis and cardiac dysrhythmias and hemorrhagic fever with coagulopathy. Chorioretinitis may also be more common than previously thought. - Skin manifestations, specifically rashes, are common; however, there are few detailed descriptions in case reports, and few images are available. Punctate erythematous, macular, and papular eruptions, most pronounced on the extremities have been observed in WNV cases and in some cases histopathologic findings have shown a sparse superficial perivascular lymphocytic infiltrate, a manifestation commonly seen in viral exanthems. A literature review provides support that this punctate rash is a common cutaneous presentation of WNV infection. WNV is one of the Japanese encephalitis antigenic serocomplex of viruses. Image reconstructions and cryoelectron microscopy reveal a 45–50 nm virion covered with a relatively smooth protein surface; this structure is similar to the dengue fever virus; both belong to the genus Flavivirus within the family Flaviviridae. The genetic material of WNV is a positive-sense, single strand of RNA, which is between 11,000 and 12,000 nucleotides long; these genes encode seven nonstructural proteins and three structural proteins; the RNA strand is held within a nucleocapsid formed from 12-kDa protein blocks; the capsid is contained within a host-derived membrane altered by two viral membrane proteins. West Nile virus has been seen to replicate faster and spread more easily to birds at higher temperatures; one of several ways climate change could impact the epidemiology of this disease. The prime method of spread of the West Nile virus (WNV) is the female mosquito. In Europe, cats were identified as being hosts for West Nile virus; the important mosquito vectors vary according to area; in the United States, Culex pipiens (Eastern United States, and urban and residential areas of the United States north of 36–39°N), Culex tarsalis (Midwest and West), and Culex quinquefasciatus (Southeast) are the main vector species. The mosquito species that are most frequently infected with WNV feed primarily on birds. Different species of mosquitos take a blood meal from different types of vertebrate hosts, Mosquitoes show further selectivity, exhibiting preference for different species of birds. In the United States, WNV mosquito vectors feed preferentially on members of the Corvidae and thrush family. Among the preferred species within these families are the American crow, a corvid, and the American robin (Turdus migratorius). Some species of birds develop sufficient viral levels (>~104.2 log PFU/ml;) after being infected to transmit the infection to biting mosquitoes that in turn go on to infect other birds. In birds that die from WNV, death usually occurs after 4 to 6 days. In mammals and several species of birds, the virus does not multiply as readily and so does not develop high viremia during infection. Mosquitoes biting such hosts are not believed to ingest sufficient virus to become infected, making them so-called dead-end hosts; as a result of the differential infectiousness of hosts, the feeding patterns of mosquitoes play an important role in WNV transmission, and they are partly genetically controlled, even within a species. Direct human-to-human transmission initially was believed to be caused only by occupational exposure, such as in a laboratory setting, or conjunctive exposure to infected blood; the US outbreak identified additional transmission methods through blood transfusion, organ transplant, intrauterine exposure, and breast feeding. Since 2003, blood banks in the United States routinely screen for the virus among their donors; as a precautionary measure, the UK's National Blood Service initially ran a test for this disease in donors who donate within 28 days of a visit to the United States, Canada, or the northeastern provinces of Italy, and the Scottish National Blood Transfusion Service asks prospective donors to wait 28 days after returning from North America or the northeastern provinces of Italy before donating. There also have been reports of possible transmission of the virus from mother to child during pregnancy or breastfeeding or exposure to the virus in a lab, but these are rare cases and not conclusively confirmed. Recently, the potential for mosquito saliva to affect the course of WNV disease was demonstrated. Mosquitoes inoculate their saliva into the skin while obtaining blood. Mosquito saliva is a pharmacological cocktail of secreted molecules, principally proteins, that can affect vascular constriction, blood coagulation, platelet aggregation, inflammation, and immunity, it clearly alters the immune response in a manner that may be advantageous to a virus. Studies have shown it can specifically modulate the immune response during early virus infection, and mosquito feeding can exacerbate WNV infection, leading to higher viremia and more severe forms of disease. Vertical transmission, the transmission of a viral or bacterial disease from the female of the species to her offspring, has been observed in various West Nile virus studies, amongst different species of mosquitoes in both the laboratory and in nature. Mosquito progeny infected vertically in autumn, may potentially serve as a mechanism for WNV to overwinter and initiate enzootic horizontal transmission the following spring, although it likely plays little role in transmission in the summer and fall. Risk factors independently associated with developing a clinical infection with WNV include a suppressed immune system and a patient history of organ transplantation. For neuroinvasive disease the additional risk factors include older age (>50+), male sex, hypertension, and diabetes mellitus. A genetic factor also appears to increase susceptibility to West Nile disease. A mutation of the gene CCR5 gives some protection against HIV but leads to more serious complications of WNV infection. Carriers of two mutated copies of CCR5 made up 4.0 to 4.5% of a sample of West Nile disease sufferers, while the incidence of the gene in the general population is only 1.0%. The most at risk occupations in the U.S. are outdoor workers, for example farmers, loggers, landscapers/groundskeepers, construction workers, painters, summer camp workers and pavers. Two reports of accidental exposure by laboratory personnel working with infected fluids or tissues have been received. While this appears to be a rare occurrence, it highlights the need for proper handling of infected materials; the World Health Organization states that there are no known cases of health care workers acquiring the virus from infected patients when the appropriate infection control precautions are observed. Preliminary diagnosis is often based on the patient's clinical symptoms, places and dates of travel (if patient is from a nonendemic country or area), activities, and epidemiologic history of the location where infection occurred. A recent history of mosquito bites and an acute febrile illness associated with neurologic signs and symptoms should cause clinical suspicion of WNV. Diagnosis of West Nile virus infections is generally accomplished by serologic testing of blood serum or cerebrospinal fluid (CSF), which is obtained via a lumbar puncture. Initial screening could be done using the ELISA technique detecting immunoglobulins in the sera of the tested individuals. Definitive diagnosis of WNV is obtained through detection of virus-specific antibody IgM and neutralizing antibodies. Cases of West Nile virus meningitis and encephalitis that have been serologically confirmed produce similar degrees of CSF pleocytosis and are often associated with substantial CSF neutrophilia. Specimens collected within eight days following onset of illness may not test positive for West Nile IgM, and testing should be repeated. A positive test for West Nile IgG in the absence of a positive West Nile IgM is indicative of a previous flavivirus infection and is not by itself evidence of an acute West Nile virus infection. If cases of suspected West Nile virus infection, sera should be collected on both the acute and convalescent phases of the illness. Convalescent specimens should be collected 2–3 weeks after acute specimens. It is common in serologic testing for cross-reactions to occur among flaviviruses such as dengue virus (DENV) and tick-borne encephalitis virus; this necessitates caution when evaluating serologic results of flaviviral infections. Four FDA-cleared WNV IgM ELISA kits are commercially available from different manufacturers in the U.S., each of these kits is indicated for use on serum to aid in the presumptive laboratory diagnosis of WNV infection in patients with clinical symptoms of meningitis or encephalitis. Positive WNV test results obtained via use of these kits should be confirmed by additional testing at a state health department laboratory or CDC. In fatal cases, nucleic acid amplification, histopathology with immunohistochemistry, and virus culture of autopsy tissues can also be useful. Only a few state laboratories or other specialized laboratories, including those at CDC, are capable of doing this specialized testing. A number of various diseases may present with symptoms similar to those caused by a clinical West Nile virus infection; those causing neuroinvasive disease symptoms include the enterovirus infection and bacterial meningitis. Accounting for differential diagnoses is a crucial step in the definitive diagnosis of WNV infection. Consideration of a differential diagnosis is required when a patient presents with unexplained febrile illness, extreme headache, encephalitis or meningitis. Diagnostic and serologic laboratory testing using polymerase chain reaction (PCR) testing and viral culture of CSF to identify the specific pathogen causing the symptoms, is the only currently available means of differentiating between causes of encephalitis and meningitis. Public health measures include taking steps to reduce mosquito populations. Personal recommendations are to reduce the likelihood of being bitten. General measures to avoid bites include: - Using insect repellent on exposed skin to repel mosquitoes. Repellents include products containing DEET and picaridin. DEET concentrations of 30% to 50% are effective for several hours. Picaridin, available at 7% and 15% concentrations, needs more frequent application. DEET formulations as high as 30% are recommended for children over two months of age; the CDC also recommends the use of: IR3535, oil of lemon eucalyptus, para-menthane-diol, or 2-undecanone. Protect infants less than two months of age by using a carrier draped with mosquito netting with an elastic edge for a tight fit. - When using sunscreen, apply sunscreen first and then repellent. Repellent should be washed off at the end of the day before going to bed. - Wear long-sleeve shirts, which should be tucked in, long trousers, socks, and hats to cover exposed skin (although most fabrics do not totally protect against bites). Insect repellents should be applied over top of protective clothing for greater protection. Do not apply insect repellents underneath clothing. - Repellents containing permethrin (e.g., Permanone) or other insect repellents may be applied to clothing, shoes, tents, mosquito nets, and other gear. (Permethrin is not suitable for use directly on skin.) Most repellent is generally removed from clothing and gear by a single washing, but permethrin-treated clothing is effective for up to five washings. - Most mosquitoes that transmit disease are most active at dawn and in the evening dusk. A notable exception is the Asian tiger mosquito, which is a daytime feeder and is more apt to be found in, or on the periphery of, shaded areas with heavy vegetation, they are now widespread in the United States, and in Florida they have been found in all 67 counties. - In an at-risk area, staying in air-conditioned or well-screened room, or sleeping under an insecticide-treated bed net is recommended. Bed nets should be tucked under mattresses, and can be sprayed with a repellent if not already treated with an insecticide. Monitoring and control West Nile virus can be sampled from the environment by the pooling of trapped mosquitoes via ovitraps, carbon dioxide-baited light traps, and gravid traps, testing blood samples drawn from wild birds, dogs, and sentinel monkeys, and testing brains of dead birds found by various animal control agencies and the public. Testing of the mosquito samples requires the use of reverse-transcriptase PCR (RT-PCR) to directly amplify and show the presence of virus in the submitted samples; when using the blood sera of wild birds and sentinel chickens, samples must be tested for the presence of WNV antibodies by use of immunohistochemistry (IHC) or enzyme-linked immunosorbent assay (ELISA). Dead birds, after necropsy, or their oral swab samples collected on specific RNA-preserving filter paper card, can have their virus presence tested by either RT-PCR or IHC, where virus shows up as brown-stained tissue because of a substrate-enzyme reaction. West Nile control is achieved through mosquito control, by elimination of mosquito breeding sites such as abandoned pools, applying larvacide to active breeding areas, and targeting the adult population via lethal ovitraps and aerial spraying of pesticides. Environmentalists have condemned attempts to control the transmitting mosquitoes by spraying pesticide, saying the detrimental health effects of spraying outweigh the relatively few lives that may be saved, and more environmentally friendly ways of controlling mosquitoes are available, they also question the effectiveness of insecticide spraying, as they believe mosquitoes that are resting or flying above the level of spraying will not be killed; the most common vector in the northeastern United States, Culex pipiens, is a canopy feeder. Eggs of permanent water mosquitoes can hatch, and the larvae survive, in only a few ounces of water. Less than half the amount that may collect in a discarded coffee cup. Floodwater species lay their eggs on wet soil or other moist surfaces. Hatch time is variable for both types; under favorable circumstances (such as warm weather), the eggs of some species may hatch in as few as 1–3 days after being laid. Used tires often hold stagnant water and are a breeding ground for many species of mosquitoes; some species such as the Asian tiger mosquito prefer manmade containers, such as tires, in which to lay their eggs. The rapid spread of this aggressive daytime feeding species beyond their native range has been attributed to the used tire trade. .No specific treatment is available for WNV infection. Most people recover without treatment. In mild cases, over-the-counter pain relievers can help ease mild headaches and muscle aches in adults. In severe cases supportive care is provided, often in hospital, with intravenous fluids, pain medication, respiratory support, and prevention of secondary infections. While the general prognosis is favorable, current studies indicate that West Nile Fever can often be more severe than previously recognized, with studies of various recent outbreaks indicating that it may take as long as 60 to 90 days to recover. Patients with milder WNF are just as likely as those with more severe manifestations of neuroinvasive disease to experience multiple somatic complaints such as tremor, and dysfunction in motor skills and executive functions for over a year. People with milder symptoms are just as likely as people with more severe symptoms to experience adverse outcomes. Recovery is marked by a long convalescence with fatigue. One study found that neuroinvasive WNV infection was associated with an increased risk for subsequent kidney disease. WNV was first isolated from a feverish 37-year-old woman at Omogo in the West Nile District of Uganda in 1937 during research on yellow fever virus. A series of serosurveys in 1939 in central Africa found anti-WNV positive results ranging from 1.4% (Congo) to 46.4% (White Nile region, Sudan). It was subsequently identified in Egypt (1942) and India (1953), a 1950 serosurvey in Egypt found 90% of those over 40 years in age had WNV antibodies; the ecology was characterized in 1953 with studies in Egypt and Israel. The virus became recognized as a cause of severe human meningoencephalitis in elderly patients during an outbreak in Israel in 1957; the disease was first noted in horses in Egypt and France in the early 1960s and found to be widespread in southern Europe, southwest Asia and Australia. The first appearance of WNV in the Western Hemisphere was in 1999 with encephalitis reported in humans, dogs, cats, and horses, and the subsequent spread in the United States may be an important milestone in the evolving history of this virus; the American outbreak began in College Point, Queens in New York City and was later spread to the neighboring states of New Jersey and Connecticut. The virus is believed to have entered in an infected bird or mosquito, although there is no clear evidence. West Nile virus is now endemic in Africa, Europe, the Middle East, west and central Asia, Oceania (subtype Kunjin), and most recently, North America and is spreading into Central and South America. Outbreaks of West Nile virus encephalitis in humans have occurred in Algeria (1994), Romania (1996 to 1997), the Czech Republic (1997), Congo (1998), Russia (1999), the United States (1999 to 2009), Canada (1999–2007), Israel (2000) and Greece (2010). Outdoor workers (including biological fieldworkers, construction workers, farmers, landscapers, and painters), healthcare personnel, and laboratory personnel who perform necropsies on animals are at risk of contracting WNV. Drought has been associated with higher number of West Nile virus causes in the following year; as drought can decrease fish and other populations that eat mosquito eggs, higher levels of mosquitoes can persist. High temperatures might also increase the risk. Higher temperatures are linked to decreased time for replication and increased viral load in birds and mosquitoes. A vaccine for horses (ATCvet code: QI05AA10 (WHO)) based on killed viruses exists; some zoos have given this vaccine to their birds, although its effectiveness is unknown. Dogs and cats show few if any signs of infection. 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School of Environment and Natural Resources 2021 Coffey Road, Columbus, Ohio 43210 Negative Effects of Livestock Grazing Riparian Areas \|James J. Hoorman OSU Extension Center at Lima Agriculture and Natural Resources \|Photo courtesy of USDA Natural Resources Conservation Service. Overgrazed sheep pasture causing gully erosion.\| The current environmental focus on controlling nonpoint pollution to protect our surface water has led to the discussion of management of riparian areas. The Environmental Protection Agency states that agriculture has a greater impact on stream and river contamination than any other nonpoint source. Grazing, particularly improper grazing of riparian areas can contribute to nonpoint source pollution. Negative impacts downstream include the contamination of drinking water supplies (55% of Ohio�s drinking water comes from surface water (Brown, 1994)), eutrophication of Lake Erie (Richards et al., 2002), and hypoxia in the Gulf of Mexico (Rabalais et al., 2001). This series of fact sheets looks at the issues of livestock and streams and what livestock producers can do to protect this precious resource. Before we discuss managing grazing livestock to decrease nonpoint pollution, it would be helpful to review the damage livestock can do to riparian areas and surface water. One cannot discuss the effects on streams by grazing livestock without recognizing the interwoven and connected nature of watersheds, riparian zones, streams, and watershed activities. Activities affecting watersheds or riparian zones also affect stream ecosystems directly, indirectly, and cumulatively. Although this series of fact sheets primarily focuses on the riparian areas, it is understated that mismanagement of the land resources in the watershed can have as big an impact on surface water. What Impact does Vegetation Removal Have on Riparian Areas? \|Photo courtesy of Jim Hoorman. Hogs in a dry creek with little vegetation and streambank stability.\| Riparian areas are the green vegetated areas adjacent to a creek, stream, or river. Riparian areas include streams, streambanks, and wetlands adjacent to streams. Impacts of vegetation removal can be placed into two categories: shifts in the plant community structure and removal of plant growth or biomass. Livestock can do both of these. Major changes in the plant community structure and usually a reduction in the number of species have been reported in the western United States. Similar or possibly even more drastic results are possible for the more humid and wet eastern United States riparian areas (Belsky, et al., 1999). Riparian vegetation has a major influence on channel shape. Vegetation increases stream bank strength by binding the soil with roots and shields banks from erosion during high flows and flooding. Kauffman and Krueger (1984) report that bank sloughing increases when vegetation removal exceeds 60%. Streams with heavily vegetated riparian areas are narrower and deeper than those that flow through poorly vegetated areas (Kauffman and Krueger, 1984). Small streams in New Zealand intensively grazed by cattle had greatly reduced shading by riparian vegetation, resulting in substantial increases in daily maximum temperatures during summer (Quinn, et al., 1992). Higher daily thermal fluctuations have also been associated with increased solar activity on the stream surface (Kauffman and Krueger, 1984). In general, differences in physical habitat and invertebrate communities were minor between paired grazed and riparian-protected reaches of the larger riparian zone, where grazing by cattle and/or sheep had little or no effect on stream shading (Quinn et al., 1992). In uplands, vegetation removal exposes soil to the energy of raindrops, facilitates sheet flow erosion with an increase in the amount of runoff and the ability to move sediment. Runoff from a heavily grazed watershed was 1.4 times higher than a moderately grazed watershed and 9 times higher than a lightly grazed watershed (Rauzi and Hanson, 1966). Summary of Effects of Vegetation Removal - Vegetation removal exposes soil to the energy of raindrops, facilitates sheet flow erosion, runoff, and the ability to move sediment. - In contrast, vegetation increases stream bank strength to resist erosion. - Stream channels along heavily vegetated areas are deeper and narrower than along poorly vegetated areas. - Sediment runoff is higher for heavily grazed watersheds compared to lightly grazed watersheds. What Impact Does Vegetation Removal Have on Water Temperature? Vegetation removal leads to higher stream water temperatures (Li et al., 1994). Riparian forest clearing in the northeastern United States resulted in increases in temperature of from 3.6 to 9.0 degrees F (Sweeney, 1993). The ungrazed stream was warmest in winter, coolest in summer, and had the narrowest range of mean daily temperature. Temperatures during summer and winter were significantly different among three streams in Pennsylvania, at least in part, related to the absence of shading due to a nearly complete lack of woody vegetation along two streams which were grazed (Wohl and Carline, 1996). Dissolved oxygen levels decline due to higher water temperatures. Algal blooms deplete oxygen by respiration at night or high oxygen demand for decomposition of algae and fecal material. This lowered oxygen environment means insufficient oxygen in spawning gravels, reduced rate of food consumption, growth, and survival of salmonids and other aquatic species, especially at their early life stages (Belsky et al, 1999). For example, it has been reported that watersheds in eastern Oregon with greater riparian canopy had higher numbers of rainbow trout (Li et al., 1994). Summary of Temperature Effects - Removal of streamside vegetation can increase mean temperature and temperature extremes. - Streams along wooded riparian zones may be cooler in summer and warmer in winter. - Relatively small changes in stream temperature can shift aquatic communities�a 3.6 degree F increase is sufficient to shift from a coldwater to a warmwater habitat. - An increase in stream temperature from 3.6 to 9 degrees F is common when streamside vegetation is removed. What Effects Does Sediment Have on Riparian Areas? \|Photo courtesy of USDA Natural Resources Conservation Service. Sedimentation from soil erosion.\| Sedimentation is recognized as the most prevalent and damaging pollution in streams in North America (Waters, 1995). Livestock grazing riparian areas can increase sediment load from the watershed, increase instream trampling, increase disturbance and erosion from overgrazed streambanks, reduced sediment trapping by riparian and instream vegetation, decreased bank stability and increased peak flows from compaction. In streams assessed in 2000, the most common agricultural pollutant was silt, which was a contributing factor for 31% of streams considered impaired (USEPA, 2000). What Impact Does Sediment Have on the Habitat for Aquatic Organisms? Sediment associated with livestock grazing occurs during snowmelt or heavy rainfall, when removal of vegetation and compaction combine to facilitate overland flow (Gardner, 1950; Bryant et al., 1972; Owens et al., 1983; Orodho et al., 1990). \|Photo courtesy of USDA Natural Resources Conservation Service. Stream with cobble and gravel bottom.\|\|Photo courtesy of USDA Natural Resources Conservation Service. Fish migrating to spawning area.\| Fine sediments increase in pools and quiet water areas from the increased erosion. Many invertebrates (insects) in streams require relatively silt-free habitats. These organisms live in the spaces between rocks in the bottom of streams (Minshall, 1984). Sediments cover and fill rocky substrates, entomb eggs and larval fish, and hinder emergence of hatched fish. Water flow in gravel is impaired, developing embryos do not receive sufficient oxygen, and metabolic wastes are not flushed. Siltation of cobble and gravel also covers hard substrates required for algal growth. This means that invertebrates that scrape algae from gravel and cobble for food will decline. It also means that invertebrates that filter food from the water column will increase. Generally, invertebrates that dwell in rock spaces are the most important food for fish that feed on invertebrates; when these species decline, so do desirable fish populations (Waters, 1995). Siltation can reduce the foraging success of aquatic organisms, fish migration can be disrupted, and respiratory systems and gills of invertebrates and fish can be impaired. Species composition and numbers of invertebrates are changed by increased sedimentation and resultant habitat changes. Pools can be filled, dam and reservoir capacity reduced, and filtration costs for domestic water supplies increased (Belsky et al., 1999). What Effects Do Livestock Grazing Have on Sediment? Sediment yield in a grazed watershed was 20-fold higher when compared to an ungrazed watershed (White et al., 1983). Sediment associated with livestock grazing occurs during snowmelt or heavy rainfall, when removal of vegetation and compaction combine to facilitate water flow into the stream (Gardner, 1950; Bryant et al., 1972; Owens et al., 1983; Orodho et al., 1990). Sedimentation from livestock grazing can be heavy enough to blanket stream beds with silt, but more commonly, leads to a gradual decrease in the depth of pools (Quinn et al., 1992; Sidle and Sharma, 1996). One comparison for sediment delivery from rotational grazing, continuous grazing, and croplands watersheds was done in Oklahoma (Olness et al., 1975). Precipitation was similar across the watersheds. Runoff ranged from 4 to 13 inches, with the highest values for the continuously grazed watersheds. Sediment delivery was highest for the continuously grazed watersheds, 8 and 10 tons of sediment per acre, and had the highest erosion index. None of the other watersheds yielded more than 4.4 tons per acre. One rotationally grazed watershed yielded the lowest sediment, 0.5 ton per acre, and the other was similar to wheat and alfalfa fields, about 1.0 ton per acre. After excluding the grazing of the banks of most perennial streams, erosion-prone hills, and pockets of native forest; sediment loads dropped by 85% in a New Zealand study (Williamson et al., 1996). A riparian zone with a diversity of vegetation is able to trap 80% to 90% of sediments transported from fields (Naiman and Decamps, 1997). For the same percent vegetative cover, more soil loss occurred from plots on steep rather than gentle slopes, and the gentle slopes could withstand more grazing pressure without seriously affecting the plant re-growth compared to steeper slopes. Slopes exceeding 5.8% are likely to suffer soil erosion even under moderate grazing pressure (Mwendera and Saleem, 1997a; Mwendera, et al., 1997). The greatest risk of summer runoff, and thus sediment yield, appears to occur in August (Owens, et al., 1989; Naeth and Chanasyk, 1996). In Ohio, annual sediment concentration decreased by more than 50% and the amount of soil lost decreased by 40% during a five-year period when cattle were fenced out of the stream relative to a seven-year period where a beef cowherd had access to a 64-acre watershed in Coshocton, Ohio (Owens et al., 1996). Average annual soil losses were reduced from 1.1 to 0.62 ton per acre while annual precipitation averages were similar during each management period. What Effect on Sedimentation Does Winter Feeding Near Riparian Areas Have? \|Photo courtesy of USDA Natural Resources Conservation Service. Winter feeding on pasture causes the highest soil erosion and highest nutrient losses.\| Winter-feeding caused a high degree of soil and plant cover disturbance and an increase in surface runoff and erosion as compared with the pastures grazed only in the summer (Chichester, et al., 1979). Feeding cattle in a winter-feeding area increased runoff and caused more chemical movement, for example, total nitrogen, total phosphorous, and organic carbon, as compared with the pastures only grazed in the summer. Evidence also suggests that cattle wintering areas may cause other related water quality problems. Winter feeding areas have shown increases in nutrients and soluble salts that can lead to development of problems with color, taste, odor, and biochemical oxygen demand (BOD). These areas can also directly produce odors from decaying products and high bacteria and pathogen loadings from animal waste. In one small, pastured watershed in eastern Ohio, runoff and sediment losses were studied for 20 years (Owens et al., 1997). In Period 1, a beef cow herd was rotationally grazed during the growing season for 12 years and was fed hay in this watershed during the dormant season. During the next three years of this study, Period 2, there was only summer rotational grazing. There was no animal occupancy on this watershed during the last five years, Period 3. Annual runoff was more than 10% of precipitation during Period 1 (4.7 inches) and less than 2% during Periods 2 and 3 (0.55 and 0.24 inches, respectively). The decrease in annual sediment loss was even greater with the change in management, yielding 2015, 130, and 8 lb per acre for the three respective periods. Over 60% of the soil loss during Period 1 occurred during the dormant (winter) season. Low amounts of runoff and erosion from three adjacent watersheds with summer-only grazing supported the conclusion that the increased runoff and erosion during Period 1 resulted from the non-rotational, continuous winter-feeding on pastures. When the management was changed, the impacts of the previous treatment were not long lasting, changing within a year. \|Photo courtesy of Jim Hoorman. Horses in a permanent pasture with continuous grazing causing streambank erosion and sedimentation and changes to stream morphology.\| In a related study in Ohio, the largest monthly average sediment concentrations were 0.8 grams per liter for 2 years without the presence of livestock. It was one and a half times higher for 3 years with 17 cows and their calves grazing during the summer months only. Sediment concentrations were four times higher for an additional six-year period with all-year grazing and hay being brought in for winter feed. Annual sediment losses were 0.09, 0.53, and 0.94 ton per acre, respectively, across the three grazing levels (Owens, et al., 1989). Summary of Sediment Effects - Sediment yield increases with increasing grazing pressure with lower levels related to ungrazed or �retired� riparian areas. - Sediment yield increases with heavy or continuous grazing, especially during the dormant season. - Soil loss increases with steeper slopes. - Sediment in streams reduces habitat for sensitive macro-invertebrates and other aquatic life. Can Grazing Livestock Affect Stream Morphology? Livestock grazing, as well as other land uses, can affect stream morphology. Stream morphology is the study of a stream�s form, structure, and channelization. There are a large number of complex, interrelated factors that determine riparian form and function that can be affected by livestock grazing which include stream discharge, sediment load, resistance of the banks and bed to movement of flowing water, vegetation, and temperature. Changes in these variables will cause an adjustment of the dynamic equilibrium of streams. What Are the Changes in Stream Morphology Due to Livestock Grazing? Streambank degradation is related both to the number of livestock grazed and the duration of grazing (Bohn and Buckhouse, 1986). Unstable stream channels and the loss of fish and invertebrate habitat are often attributed to cattle grazing practices in riparian areas in the western United States. Cattle grazing often cause large changes in channel morphology, causing wider, shallower stream channels (Knapp et al., 1998) with significant native vegetation overhang and extensive fish habitat changing to wide braided channels with little cover for fish or amphibians (Williamson et al., 1992). A number of studies have examined effects of livestock by fencing riparian areas to exclude grazing and then noting the effects on riparian vegetation and stream morphology. Magilligan and McDowell (1997) selected four gravelbedded, steep streams in eastern Oregon and excluded cattle for 14 years to study stream changes. Reductions in bankfull widths by 10 to 20 percent and increases of 8 to 15 percent in pool area were the most common and identifiable changes in excluding the cattle. Not all channel properties demonstrated adjustment, leading Magilligan and McDowell (1997) to suggest that perhaps 14 years is an insufficient duration for these variables to adjust. \|Photo courtesy of USDA Natural Resources Conservation Service. Pools and riffles create fish habitat.\| In an examination of two grazed and one ungrazed reaches where grazing had been excluded for 11 years, the ungrazed section demonstrated improved riparian vegetation and a deeper, narrower channel. However, exclusion of grazing for two years along a 0.6 mile section of a stream channel that had experienced historical grazing did not lead to substantial stream recovery, indicating that it may take many years of excluding livestock from streams for the full benefits to be realized by the stream (Platts and Nelson, 1985). Shifting the location of cattle grazing can cause significant downstream impacts. Cattle moved into wet riparian areas upstream caused decreases in thalweg depth (i.e., the location of the deepest portion of the channel), increases in fine sediment deposition in the channel, and loss of pool volume in these upstream areas. It was also reported that the deposition of fine sediment in reaches with high volumes of large woody debris increased (Sidle and Sharma, 1996). Pool/riffle ratio (a measure of fish habitat), as well as soil and vegetation stability, varied significantly with cattle density (Meyers and Swanson, 1991). As pool/riffle ratios changed, other channel properties are seen to change as well. Continuous, heavy grazing resulted in a stream reach (low part of the stream bank adjacent to a stream) that became four times wider and one-fifth as deep as an adjacent area that was only lightly grazed (Platts and Wagstaff, 1984). What Are the Stream Stability Changes Due to Grazing? The removal of riparian vegetation has severe effects on stream channel characteristics. Streambank stability is reduced due to fewer plant roots to anchor soil, less plant cover to protect the soil surface from erosional disturbance and the shear force of trampling hooves. Impacts include increased streambank sloughing, increased erosion, increased channel width, and reduced depth. Streambank undercuts are reduced due to streambank breakdown by sloughing and trampling. The stream channel contains fewer meanders and gravel bars due to increased water velocity. Pools decrease in number and quality from increased sediment and loss of woody debris (Belsky et al., 1999). Comparisons of grazed and ungrazed streams found that grazed stream channels tend to be wider with shallower banks (Marcuson, 1977; Duff, 1979). When animals graze directly on streambanks, mass erosion from trampling, hoof slide, and streambank collapse causes soil to move directly into the stream (Platts, 1991). \|Photo courtesy of USDA Natural Resources Conservation Service. Meandering streams create fish habitat.\| \|Photo courtesy of Jim Hoorman. Horses in the stream deposit urine and feces causing excess nutrients and eutrophication. Manure deposited along or directly into streams elevates concentrations of phosphorus and nitrogen (Lemly, 1982).\| The loss of stream channel integrity and diversity results in impacts to fish populations. For example, Marcuson (1977) studied the difference in habitat and fish populations in grazed and ungrazed stream sections. The study documented 80% more stream alteration in the grazed area than in an adjacent ungrazed area with the grazed area losing 11 acres of a 120-acre pasture. The ungrazed section produced 256 more pounds of fish per acre than the grazed section. Summary of Effects Due to Channel Morphology - Unstable stream channels and the loss of fish and invertebrate habitat are often attributed to cattle grazing practices in riparian areas. - Stream channels along heavily vegetated areas are deeper and narrower than along poorly vegetated areas. - Livestock management often causes local changes in habitat, thereby impacting fish and invertebrates. - Changes are much more pronounced in small streams than large ones; impacts on lakes are under-studied but appear to be minimal. - The natural variance among stream channels, lakes, and wetlands makes generic conclusions very difficult. Most impacts and most Best Management Practices will be site-specific. Site-specific BMPs depend on stream morphology. What Impact Does Eutrophication Have on Riparian Areas? The presence of some aquatic vegetation is normal in streams and indicates a healthy stream. Algae and aquatic plants provide habitat and food for all stream animals. High levels of nutrients (especially phosphorus and nitrogen) promote an overabundance of algae and floating and rooted aquatic plants. An excessive amount of aquatic vegetation is not beneficial to most stream life. Plant respiration and decomposition of dead plant life consume dissolved oxygen in the water. Eutrophication is the process where aquatic vegetation grows quickly and decomposes, consuming oxygen from the stream. Lack of dissolved oxygen creates stress for all aquatic organisms and can cause fish kills. Elevated levels of nutrient input often results in dense growths of filamentous green algae, i.e., Cladophora spp. These dense growths promote the production of some insect species, and replace diverse populations of attached plant microorganisms. Many herbivorous insects decline greatly in response to dense algae blooms (Li et al., 1994). The dense filamentous algae reduce feeding efficiency of insect eating fish due to a switch in the insect prey base (Tait et al., 1994). What Are the Nutrient Loads Associated with Grazing? Concentrations of Ammonium Nitrogen (NH4-N), total Kjeldahl N, and total P, were directly related to the density of grazing livestock. Leachates from the standing plant material, surface litter layer, surface soil, and manure deposits indicated manure and standing plant material were likely sources of most chemical components in runoff water (Schepers et al., 1982). Over-application of fertilizer and manure can overload the soil with phosphorous. Iron, aluminum, and calcium in the soil bind excess phosphorous. In flooded soils, iron binds less phosphorous than it does in drier, aerobic soils. This decreased binding ability increases the availability of phosphorous for plant uptake and for movement into surface water (Green and Kaufman, 1989). Since riparian areas have limited ability to hold excess phosphorous, they are relatively ineffective in protecting streams against poor phosphorous management practices on upland areas. Thus, good upland management is necessary to protect against phosphorous pollution (Bellow, 2003). Owens et al. (1994) found that the various forms of nitrogen (N) increased in ground water during a five-year period with 200 lb N per acre annual fertilizer application to a grass-pasture grazed by beef cattle and reached levels that were usually in excess of 10 PPM. Ohio EPA has set 10 PPM nitrate-nitrogen as the upper threshold for drinking water. Nitrate concentrations in groundwater dropped rapidly after alfalfa was inter-seeded into the grass pastures and N fertilizer was no longer applied. The amount of N lost via subsurface flow decreased, but subsurface flow remained the main pathway for N loss compared with surface runoff or sediment-attached N. In a related study in Ohio, Owens et al. (1989) found nutrient concentrations remained low across three grazing levels, with the exception of potassium (K) concentration, which increased with all-year grazing. \|Photo courtesy of USDA Natural Resources Conservation Service. The uneven recycling of N through feces and urine may increase nitrate leaching. The extent to which nitrate can leach from beneath urine and fecal spots under soil and climatic has not been studied extensively.\|\|Photo courtesy of USDA Natural Resources Conservation Service. Cattle in feedlot with stream causing soil erosion and eutrophication.\| The nitrate losses in a N-fertilized orchardgrass field in central Pennsylvania averaged across three years were 10.4 lb per acre for the control, 15.0 lb per acre for manure applied, 196.2 lb per acre for urine applied in spring, 214.1 lb per acre for summer applied urine, and 281.0 lb per acre for fall-applied urine (Stout, et al., 1997). These losses represent about 2% of the N applied in the feces and about 18%, 28%, and 31% of the spring-, summer-, and fall-applied urine N. Winter-feeding caused a high degree of soil and plant cover disturbance and an increase in surface runoff and erosion as compared with the pastures grazed only in the summer. Feeding cattle in a winter-feeding area increased runoff that caused more chemical movement of total N, total P, and organic carbon as compared with the pastures only grazed in the summer (Chichester et al., 1979). The level of nitrogen leaching as the result of urination by cows depended upon soil type, moisture conditions, and grazing intensity in the Netherlands. Nitrate levels rarely exceeded 11.3 PPM in a wet moderately grazed field, but were often exceeded in a drier more heavily grazed pasture (Hack-ten Broeke et al., 1996). Total nitrogen in runoff from a cattle-grazed watershed in the Pacific northwest ranged between < 0.9 to 3.6 lb per acre per year, whereas only 0.44 lb per acre per year was measured in a non-grazed watershed over a three-year period (Jawson et al., 1982). Total nitrogen received in precipitation was equal to or greater than nitrogen lost in runoff from the grazed watershed. Nitrate-N levels in the runoff were normally less than 1 mg/l. Total P (TP) losses in runoff from the grazed watershed ranged from 0.09-1.2 lb per acre per year and from < 0.09 to 0.15 lb per acre per year from the ungrazed area (Jawson et al., 1982). Concentrations of ammonia nitrogen (NH4-N), nitrate nitrogen (NO3-N), total phosphorous (TP), soluble P, and chloride (Cl-) in runoff from a 6.2-acre cow-calf grazed pasture in Nebraska were 6%, 45%, 37%, 48%, and 78% greater, respectively, over a three-year period when livestock were grazing in comparison to periods when cattle were removed (Schepers and Francis, 1982). Mean concentrations of Kjeldahl N, total P, Calcium (Ca), and potassium (K) were 190, 150, 24, and 240 times higher, respectively, in surface runoff from a grazed and fertilized hill pasture than those in rainfall over one year. Concentrations tended to be high in summer and were strongly related to grazing. The peak of nitrate concentration after grazing lagged approximately two weeks after the Kjeldahl N peak and probably depended on nitrification of ammonium from dung and urine (McColl and Gibson, 1979). Do Riparian Zones Impact Nutrient Flow? \|Photo courtesy of USDA Natural Resources Conservation Service. Riparian forest buffers and protects stream water quality.\| Dissolved nutrients are transported into streams primarily in the groundwater (Gregory et al., 1991). Because of the riparian zone position within the watershed, it intercepts the soil solution as it passes through the rooting zone prior to entering the stream. Riparian zones also contribute seasonal pulses of dissolved components derived from plant litter into streams. Thus, the riparian zone functions to remove nutrients and modify inputs to the stream. Riparian forests were responsible for removal of more than three-quarters of the dissolved nitrate transported from Photo courtesy of USDA Natural Resources Conservation Service. Riparian forest buffers and protects stream water quality. croplands into a Maryland river (Peterjohn and Correll, 1984). Natural riparian forests can denitrify and release 25 to 35 pounds of nitrogen per acre per year (Cole, 1981). Because of their unique position at the interface between terrestrial (land) and aquatic ecosystems, riparian zones play a critical role in controlling the flow of nutrients from watersheds. What Effect Does Livestock Exclusion Have on Riparian Areas? Sediment, phosphorus, particulate- and nitrate-nitrogen concentrations, during 71 run-off events over 22 months, were lower and varied significantly less at retired riparian pasture than at grazed riparian pasture sites (Smith, 1989). Riparian pasture retirement is an effective means of reducing surface runoff pollutant loads to waterways. After exclusion of grazing (�retirement�) from the banks of most perennial streams, erosion-prone hills, and remnant pockets of forest in a New Zealand watershed; loads decreased by 27% for particulate P, 26% for soluble P, 40% for particulate N, and increased by 26% for dissolved N (Williamson et al., 1996). Williamson et al. (1996) predicted that retirement reduced total phosphorous loads by 20% in the lake that receives runoff from the watershed. During 12 years after retirement from grazing, dominant vegetation in the set-aside areas changed from pasture grasses to native species in pasture along the edge of a small stream in New Zealand (Cooper et al., 1995). The riparian set-aside soils had higher water conductivity indicating that surface runoff water transported into the zone would infiltrate, fill soil pores, and emerge as subsurface flow at the stream edge. This research implied that riparian set-aside has led to the development of a zone likely to supply runoff to the adjacent stream that is depleted in sediment-bound nutrients and dissolved Nitrogen (N) but enriched in dissolved Phosphorous (P). Summary of Nutrient Effects - Excess nutrients in streams cause eutrophication to increase. Eutrophication is the process where aquatic vegetation grows quickly and decomposes, consuming oxygen from the stream. - Nutrient concentrations (various forms of N and P) in runoff increase with increasing grazing duration. - Retiring areas from grazing but maintaining grass vegetation reduces nutrient delivery, but dissolved N may be reduced differentially in relation to dissolved P. For more information on the effects of livestock grazing riparian areas, see the following fact sheets: - Understanding the Benefits of Healthy Riparian Areas, LS-1-05 - The Effects of Grazing Management on Riparian Areas, LS-3-05 - Best Management Practices to Control the Effects of Livestock Grazing Riparian Areas, LS-4-05 - Pathogenic Effects from Livestock Grazing Riparian Areas, LS-5-05 This fact sheet was adapted from Generic Environmental Impact Statement on Animal Agriculture: A Summary of Literature Related to the Effects of Animal Agriculture on Water Resources (G), 1999, Univ. of Minnesota. Bellows, B. C. March 2003. Protecting riparian areas: Farmland management strategies. Soil Systems Guide, Appropriate Technology Transfer for Rural Areas. At www.attra.ncat.org. Belsky, A. J., A. Matzke, and S. Uselman. 1999. Survey of livestock influences on stream and riparian ecosystems in the western United States. Journal of Soil and Water Conservation 54(1): 419-431. Bohn, C. C., and J. C. Buckhouse. 1986. Effects of grazing management on streambanks. Trans. N. Am. Wildl. Natl. Resour. Conf. 51:265-271. Bryant, H. T., R. E. Blaser, and J. R. Peterson. 1972. Effect of trampling by cattle on bluegrass yield and soil compaction of a meadowville loam. Agron. J. 64:331-334. Chichester, F. W., R. W. Van Keuran, and J. L. McGuinness. 1979. Hydrology and chemical quality of flow from small pastured watersheds: Chemical quality. J. Envir. Qual. 8(2): 167-171. Cole, D. W., 1981. Nitrogen uptake and translocation by forest ecosystems. In: F. E. Clark and T. Rosswall (eds.) Terestrial Nitrogen Cycles. Ecological Bulletin. Vol. 33. p. 219-232. Cooper, A. B., C. M. Smith, and M. J. Smith. 1995. Effects of riparian set-aside on soil characteristics in an agricultural landscape�Implications for nutrient transport and retention. Agric. Ecosystems Environ. 55:61-67. Duff, Donald A. 1979. Riparian habitat recovery on Big Creek, Rich County, Utah. In Proceedings: Forum�Grazing and Riparian/Stream Ecosystems. Trout Unlimited, Inc. p. 91 Gardner, J. L. 1950. Effects of thirty years of protection from grazing in desert grassland. Ecology. 31:44-50. Generic Environmental Impact Statement on Animal Agriculture: A Summary of Literature Related to the Effects of Animal Agriculture on Water Resources (G), 1999. The Environmental Quality Board, College of Agriculture, Food, and Environmental Sciences (COAFES), Univ. of Minnesota. Green, D. M., and J. B. Kauffman. 1989. Nutrient cycling at the land-water interface: The importance of the riparian zone. In: R. E. Gresswell, B. A. Barton, and J. L. Kershner (eds.) Practical Approaches to Riparian Resource Management: An Education Workshop. U.S. Bureau of Land Management. Billings, MT. p. 61-68. Gregory, S. V., F. J. Swanson, W. A. McKee, and K. W. Cummins. 1991. An ecosystem perspective of riparian zones. Bioscience 41(8): 540-550. Hack-ten Broeke, M. J. D., W. J. M. De Groot, and J. P. Dijkstra. 1996. Impact of excreted nitrogen by grazing cattle on nitrate leaching. Soil Use Manage. 12:190-198. Jawson, M. D., L. F. Elliott, K. E. Saxton, and D. H. Fortier. 1982. The effect of cattle grazing on nutrient losses in a pacific northwest setting, USA. J. Environ. Qual. 11:628-631. Kaufmann, J. B., and W. C. Kreuger. 1984. Livestock impacts on riparian ecosystems and streamside management implications: A review. J. Range Manage. 37:430-438. Knapp, R. A., V. T. Vredenburg, and K. R. Matthews. 1998. Effects of stream channel morphology on golden trout spawning habitat and recruitment. Ecol. Appl. 8:1104-1117. Lemly, D. A. 1982. Modification of benthic insect communities in polluted streams: Combined effects of sedimentation and nutrient enrichment. Hydrobiologia. 87:229-245. Li, H. W., G. A. Lamberti, T. N. Pearsons, C. K. Tait, J. L. Li, and J. C. Buckhouse. 1994. Cumulative effects of riparian disturbances along high desert trout streams of the John Day Basin, Oregon. Trans. Am. Fisheries Soc. 123:627-640. Magilligan, F. J., and P. F. McDowell. 1997. Stream channel adjustments following elimination of cattle grazing. J. Am. Water Resour. Assn. 33:867-878. Marcuson, Patrick E. 1977. Overgrazed streambanks depress fishery production in Rock Creek, Montana. Fish and Game Federation Aid Program. F-20-R-21-11a. McColl, R. H. S., and A. R. Gibson. 1979. Downslope movement of nutrients in hill pasture,Taita, New Zealand: 2. Effects of season, sheep grazing and fertilizer. New Zealand J. Agric. Res. 22:151-162. Meyers, T. J., and S. Swanson. 1991. Aquatic habitat condition index, streamtypes and livestock bank damage in northern Nevada. Water Resour. Bull. 27:667-677. Minshall, G. W. 1984. Aquatic insect substratum relationships. In V. H. Resh and D. M. Rosenberg (ed.) The ecology of aquatic insects. Praeger Publishers, New York. p. 356-400. Mwendera, E. J., and M. A. M. Saleem. 1997a. Infiltration rates, surface runoff, and soil loss as influenced by grazing pressure in the Ethiopian highlands. Soil Use Manage. 13:29-35. Mwendera, E. J., M. A. M. Saleem, and A. Dibabe. 1997. The effect of livestock grazing on surface runoff and soil erosion from sloping pasture lands in the Ethiopian highlands. Australian J. Experimental Agric. 37:421-430. Naeth, M. A., and D. S. Chanasyk. 1996. Runoff and sediment yield under grazing in foothills fescue grasslands of Alberta. Water Res. Bull. 32:89-95. Naiman, R. J., and H. Decamps. 1997. The ecology of interfaces: Riparian zones. Annual Review of Ecology and Systematics. V. 28. p. 621-658. Olness, A., S. J. Smith, E. D. Rhoades, and R. G. Menzel. 1975. Nutrient and sediment discharge from agricultural watersheds in Oklahoma. J. Environ. Qual. 4:331-336. Ohio�s Hydrologic Cycle. 1994. L. C. Brown. AEX 461. Ohio State University Extension. Orodho, A. B., M. J. Trlica, and C. D. Bonham. 1990. Long term heavy grazing effects on soil and vegetation in the four corners region. Southwest Naturalist. 35:9-14. Owens, L. B., W. M. Edwards, and R. W. Van Keuren. 1989. Sediment and nutrient losses from an unimproved all-year grazed watershed. J. Environ. Qual. 18:232-238. Owens, L. B., W. M. Edwards, and R. W. Van Keuren. 1996. Sediment losses from a pastured watershed before and after stream fencing. J. Soil Water Conserv. 51:90-94. Owens, L. B., W. M. Edwards, and R. W. Van Keuren. 1997. Runoff and sediment losses resulting from winter feeding on pastures. J. Soil Water Conserv. 52:194-197. Owens, L. B., W. M. Edwards, and R. W. Van Keuren. 1983. Surface runoff quality comparisons between unimproved pasture and woodlands. J. Environ. Qual. 12:518-522. Owens, L. B., W. M. Edwards, and R. W. Van Keuren. 1994. Groundwater nitrate levels under fertilized grass and grasslegumes pastures. J. Environ. Qual. 23:752-758. Richards, R. P., F. G. Calhoun, and G. Matisoff. 2002. Lake Erie agricultural systems for environmental quality project. J. of Envir. Qual. 31:6-16. Rabalais, N. N., R. E. Turner, and W. J. Wiseman, Jr. 2001. Hypoxia in Gulf of Mexico. J. of Envir. Qual. Mar-Apr 30(2):320-329. Platts, W. S. 1991. Livestock grazing. In: Influence of forest and rangeland management on Salmonid fishes and their habitats. American Fisheries Society, Special Publication 19:389-423. Platts, W. S., and R. F. Nelson. 1985. Stream habitat and fisheries response to livestock grazing and instream improvement structures, Big Creek, Utah. J. Soil Water Conserv. 40:374-379. Platts, W. S. and F. J. Wagstaff. 1984. Fencing to control livestock grazing on riparian habitats along streams: Is it a viable alternative. N. Am. J. Fisheries Manage. 4:266-272. Peterjohn, W. T., and D. L. Correll. 1984. Nutrient dynamics in an agricultural watershed: Observations of a riparian forest. Ecology 65: 1466-1475. Quinn, J. M., R. B. Williamson, R. K. Smith, and M. L. Vickers. 1992. Effects of riparian grazing and channelization on streams in southland New Zealand 2. Benthic invertebrates. New Zealand J. Marine Freshwater Res. 26:259-273. LS-2-05�page 10 Rauzi, F., and C. L. Hanson. 1966. Water intake and runoff as affected by intensity of grazing. J. Range Manage. 19:351-356. Schepers, J. S., and D. D. Francis. 1982. Chemical water quality of runoff from grazing land in Nebraska: I. Influence of grazing livestock. J. Environ. Qual. 11:351-354. Schepers, J. S., B. L. Hackes, and D. D. Francis. 1982. Chemical water quality of runoff from grazing land in Nebraska: II. Contributing factors. J. Environ. Qual. 11:355-359. Sidle, R. C., and A. Sharma. 1996. Stream channel changes associated with mining and grazing in the Great Basin. J. Environ. Qual. 25:1111-1121. Smith, C. M. 1989. Riparian pasture retirement effects on sediment phosphorus and nitrogen in channellized surface run-off from pastures. New Zealand J. Mar. Freshwater Res. 23:139-146. Stout, W. L., S. A. Fales, L. D. Muller, R. R. Schnabel, W. E. Priddy, and G. F. Elwinger. 1997. Nitrate leaching from cattle urine and feces in northeastern U.S. Soil Sci. Soc. Am. 61:1787. Sweeny, B. W. 1993. Effects of streamside vegetation on macroinvertebrate communities of White Clay Creek in eastern North America. Proc. of the Natural Science Academy of Philadelphia. 144:291-340. Tait, C. K., J. L. Li, G. A. Lamberti, T. N. Pearsons, and H. W. Li. 1994. Relationships between riparian cover and community structure of high desert streams. J. N. Am. Benthological Soc. 13:45-56. USEPA. 2000. National Water Quality Inventory: 2000 Report to Congress Executive Summary, Office of Water, Washington, DC 20460. [Online] Available at http://www.epa.gov/305b. Waters, T. F. 1995. Sediment in streams, sources, biological effects and control. American Fisheries Society Monograph 7. White, R. K., R. W. VanKeuren, L. B. Owens, W. M. Edwards, and R. H. Miller. 1983. Effects of livestock pasturing on non-point surface runoff. Project Summary, Robert S. Kerr Environmental Research Laboratory, Ada, Oklahoma. EPA- 600/S2-83-011. 6p. Williamson, R. B., C. M. Smith, and A. B. Cooper. 1996. Watershed riparian management and its benefits to a eutrophic lake. J. Water Res. Planning Manage.-ASCE. 122:24-32. Williamson, R. B., R. K. Smith, and J. M. Quinn. 1992. Effects of riparian grazing and channelization on streams in Southland New Zealand I. Channel form and stability. New Zealand Journal of Marine & Freshwater Research. 26:241-258. Wohl, N. E., and R. F. Carline. 1996. Relations among riparian grazing, sediment loads, macroinvertebrates, and fishes in three central Pennsylvania streams. Can. J. Fisheries Aquatic Sci. 53(suppl. 1):260-266. The following persons reviewed the original material: Dr. Lloyd Owens, Soil Scientist, USDA-ARS; Dr. Steve Loerch, Professor of Animal Sciences, The Ohio State University; and Robert Hendershott, Grassland Specialist, USDA-NRCS. The authors would also like to thank Jerry Iles, Extension Educator, Watershed Management, OSU Center at Piketon for comments and suggestions. The authors thank Kim Wintringham (Technical Editor, Section of Communications and Technology) for editorial and graphic production. Click here for PDF version of this Fact Sheet. OSU Extension embraces human diversity and is committed to ensuring that all educational programs conducted by Ohio State University Extension are available to clientele on a nondiscriminatory basis without regard to race, color, age, gender identity or expression, disability, religion, sexual orientation, national origin, or veteran status. Keith L. Smith, Associate Vice President for Agricultural Administration and Director, OSU Extension
This week in precalc 11, we learned many new concepts and learned a completely different vocabulary. We learned something called reciprocal functions. For example, if we have then the reciprocal would be Some new vocabulary that we learned are invariant points, hyperbola, and asymptotes Invariant point : 1 and -1 are invariant points because once reciprocated will stay the same unlike other numbers. Hyperbola : this is the graph of a reciprocal function Asymptotes : the vertical and horizontal asymptotes are the boundaries where the hyperbola can not touch or overpass Hyperbolas are quite interesting and the way to tell if they are indeed hyperbolas is to see if there are two parts to the graph. Here is an example of a hyperbola (linear) : as you can see there are two parts to the graph For this graph, we can see that the linear graph is going through a point of -3 on the x intercept. This indicates the vertical asymptotes of the hyperbola. As for the horizontal asymptote (y), in our grade the y will always equal to 0. So for our asymptotes for this graph, it would be and . To find the invariant points, on the y axis we find 1 and -1 and from there, move horizontally until we hit a point of the linear graph. Those will be our invariant points, in this case our invariant points are and . What I just showed were hyperbolas for linear graphs but there are also hyperbolas for quadratics. They are both pretty similar to each other in regards to the steps.
Socratic Seminar. Socrates (June 4, 470 BC – May 7, 399 BC) (Greek Σωκράτης Sōkrátēs; invariably anglicized as /'sɒkɹətiːz/ Sǒcratēs) was a Greek (Athenian) philosopher. Socrates (June 4, 470 BC – May 7, 399 BC) (Greek Σωκράτης Sōkrátēs; invariably anglicized as /'sɒkɹətiːz/ Sǒcratēs) was a Greek (Athenian) philosopher. The Socratic method of teaching is based on Socrates' theory that it is more important to enable students to think for themselves than to merely fill their heads with "right" answers. Therefore, he regularly engaged his pupils in dialogues by responding to their questions with questions, instead of answers. This process encourages divergent thinking rather than convergent. Students are given opportunities to "examine" a common piece of text, whether it is in the form of a novel, poem, art print, or piece of music. After "reading" the common text "like a love letter,” open-ended questions are posed. Open-ended questions allow students to think critically, analyze multiple meanings in text, and express ideas with clarity and confidence. After all, a certain degree of emotional safety is felt by participants when they understand that this format is based on dialogue and not discussion/debate. Dialogue is exploratory and involves the suspension of biases and prejudices. Discussion/debate is a transfer of information designed to win an argument and bring closure. Americans are great at discussion/debate. We do not dialogue well. However, once teachers and students learn to dialogue, they find that the ability to ask meaningful questions that stimulate thoughtful interchanges of ideas is more important than "the answer." Participants in a Socratic Seminar respond to one another with respect by carefully listening instead of interrupting. Students are encouraged to "paraphrase" essential elements of another's ideas before responding, either in support of or in disagreement. Members of the dialogue look each other in the "eyes" and use each other names. This simple act of socialization reinforces appropriate behaviors and promotes team building. For Participants in a Socratic Seminar Socrates after being sentenced to die for impiety,introducing new gods, and corrupting the young. 1. Refer to the text when needed during the discussion. A seminar is not a test of memory. You are not "learning a subject;” your goal is to understand the ideas, issues, and values reflected in the text. 2. It's OK to "pass" when asked to contribute. 3. Do not participate if you are not prepared. A seminar should not be a bull session. 4. Do not stay confused; ask for clarification. 5. Stick to the point currently under discussion; make notes about ideas you want to come back to. 6. Don't raise hands; take turns speaking. 7. Listen carefully. 11. You are responsible for the seminar, even if you don't know it or admit it. Of Participants in a Socratic Seminar "Socrates said he was not an Athenian or a Greek, but a citizen of the world." Dialogue Vs. Debate What IS the difference? Dialogue is collaborative: multiple sides work toward shared understanding. In dialogue, one listens to understand, to make meaning, and to find common ground. Dialogue enlarges and possibly changes a participant's point of view. Dialogue creates an open-minded attitude: an openness to being wrong and an openness to change. In dialogue, one submits one's best thinking, expecting that other people's reflections will help improve it rather than threaten it. Dialogue calls for temporarily suspending one's beliefs. In dialogue, one searches for strengths in all positions. Dialogue respects all the other participants and seeks not to alienate or offend. Dialogue assumes that many people have pieces of answers and that cooperation can lead to a greater understanding. Dialogue remains open-ended. Debate is oppositional: two opposing sides try to prove each other wrong. In debate, one listens to find flaws, to spot differences, and to counter arguments. Debate defends assumptions as truth. Debate creates a close-minded attitude, a determination to be right. In debate, one submits one's best thinking and defends it against challenge to show that it is right. Debate calls for investing wholeheartedly in one's beliefs. In debate, one searches for weaknesses in the other position. Debate rebuts contrary positions and may belittle or deprecate other participants. Debate assumes a single right answer that somebody already has. Debate demands a conclusion. How do I earn a grade "Wisdom begins in wonder."
Bivariate data arises when a study aims to determine whether there is a relation between two variable quantities. The quantities under investigation are called the explanatory variable and the response variable, or equivalently, the independent- and the dependent variable. If the dependent variable is found to be related in a definite way to the values taken by the independent variable, then further research may show that there is a causal relationship between the two. This is not necessarily the case, however, because the two quantities may both be varying in response to changes in a third factor and so it could not be claimed that one of the variables being studied has had a causal effect on the other. You can read more on causality and correlation here. The results of bivariate data investigations can be displayed graphically using a scatter plot. The level of the independent variable corresponds to a position on the horizontal axis and the resulting value of the dependent variable corresponds to a distance along the vertical axis. In this way, each data point is displayed as a point in a two-dimensional coordinate system. A correlation is a way of expressing a relationship between two variables and, more specifically, how strongly pairs of data are related. We describe the correlation from data using language like positive correlation, negative correlation or no correlation. We can even further strengthen the language by using strong or weak. It is often possible to determine by looking at a scatter plot the nature of a relation between variables. If the data points lie on or close to a line, a linear relation is strongly suggested. An algebraic model of the form $y=ax+b$y=ax+b may then be proposed as a summary of the relation. If the points are not close to a line but still display a generally linear trend, a weak linear relation may be said to exist. Techniques are available to find the best fitting linear model to any bivariate data set whether or not there is a genuine linear relation. The slope of a line fitted to a data set may be positive or negative, depending on the sign of the coefficient $a$a in the formula $y=ax+b$y=ax+b. This corresponds to whether the response variable increases or decreases respectively in response to an increase in the explanatory variable. A positive correlation is when the data appears to gather in a positive relationship. Similar to a straight line with a positive gradient. In other words, as one variable increases, the other variable also increases. There are three types of positive correlation: You may also come across a moderate correlation, which is a correlation between weak and strong. Here are some examples of positive style correlations A negative correlation is when the data appears to gather in a negative relationship. Similar to a straight line with a negative gradient. In other words, as one variable increases, the other one decreases. Like positive correlation, there are three types of negative correlation: Here are some examples of negative style correlations No correlation is when there is no relationship between the variables. This means that there is a random or nonlinear relationship between the two sets of data. Here is a diagram of no correlation Shapes other than a line may be apparent in a scatter plot. If the data points lie on or near a curve, it may be appropriate to infer a non-linear relation between the variables. It is possible, for example, to find the best-fitting polynomial of a given degree or some other function, that reasonably describes the observed effect. Non-linear relations would still be described as strong/moderate/weak and positive/negative depending on how strongly they resemble the chosen curve and whether the curve is positive or negative in shape. Other features that may appear in scatter plots include clustering and outliers. In an observational study, a gap in the values available in the explanatory variable may create the appearance of clusters in the response values. Such a feature might arise when the population under consideration is made up of distinct sub-populations. Even without gaps in the values of the independent variable there may be distinct sets of data points in which different trends are apparent, indicating the existence of different groups within the population. An outlier occurs where a single value of the response variable is very different from neighbouring values. An outlier might be due to measurement error, but not necessarily. An outlier should not be discarded from the data set before looking for a satisfactory explanation. Consider the following graph: The correlation is: (select the best answer) The correlation features: The scatter plot shows the relationship between sea temperatures and the amount of healthy coral. Describe the correlation between sea temperature the amount of healthy coral. Select all that apply. Which variable is the dependent variable? Which variable is the independent variable? The following table has data results from an experiment. Plot the data from the table on the graph below. What is the type of correlation between the data points? Select the best answer. Plan and conduct investigations using the statistical enquiry cycle: A justifying the variables and measures used B managing sources of variation, including through the use of random sampling C identifying and communicating features in context (trends, relationships between variables, and differences within and between distributions), using multiple displays D making informal inferences about populations from sample data E justifying findings, using displays and measures. Investigate a given multivariate data set using the statistical enquiry cycle Investigate bivariate numerical data using the statistical enquiry cycle
October 26, 2015 – Using a new process in planetary formation modeling, where planets grow from tiny bodies called “pebbles,” Southwest Research Institute scientists can explain why Mars is so much smaller than Earth. This same process also explains the rapid formation of the gas giants Jupiter and Saturn, as reported earlier this year. “This numerical simulation actually reproduces the structure of the inner solar system, with Earth, Venus, and a smaller Mars,” said Hal Levison, an Institute scientist at the SwRI Planetary Science Directorate. He is the first author of a new paper published in the Proceedings of the National Academy of Sciences of the United States (PNAS) Early Edition. The fact that Mars has only 10 percent of the mass of the Earth has been a long-standing puzzle for solar system theorists. In the standard model of planet formation, similarly sized objects accumulate and assimilate through a process called accretion; rocks incorporated other rocks, creating mountains; then mountains merged to form city-size objects, and so on. While typical accretion models generate good analogs to Earth and Venus, they predict that Mars should be of similar-size, or even larger than Earth. Additionally, these models also overestimate the overall mass of the asteroid belt. “Understanding why Mars is smaller than expected has been a major problem that has frustrated our modeling efforts for several decades,” said Levison. “Here, we have a solution that arises directly from the planet formation process itself.” New calculations by Levison and co-authors Katherine Kretke, Kevin Walsh and Bill Bottke, all of SwRI’s Planetary Science Directorate follow the growth and evolution of a system of planets. They demonstrate that the structure of the inner solar system is actually the natural outcome of a new mode of planetary growth known as Viscously Stirred Pebble Accretion (VSPA). With VSPA, dust readily grows to “pebbles” — objects a few inches in diameter — some of which gravitationally collapse to form asteroid-sized objects. Under the right conditions, these primordial asteroids can efficiently feed on the remaining pebbles, as aerodynamic drag pulls pebbles into orbit, where they spiral down and fuse with the growing planetary body. This allows certain asteroids to become planet-sized over relatively short time scales. However, these new models find that not all of the primordial asteroids are equally well-positioned to accrete pebbles and grow. For example, an object the size of Ceres (about 600 miles across), which is the largest asteroid in the asteroid belt, would have grown very quickly near the current location of the Earth. But it would not have been able to grow effectively near the current location of Mars, or beyond, because aerodynamic drag is too weak for pebble capture to occur. “This means that very few pebbles collide with objects near the current location of Mars. That provides a natural explanation for why it is so small,” said Kretke. “Similarly, even fewer hit objects in the asteroid belt, keeping its net mass small as well. The only place that growth was efficient was near the current location of Earth and Venus.” “This model has huge implications for the history of the asteroid belt,” said Bottke. Previous models have predicted that the belt originally contained a couple of Earth-masses’ worth of material, meaning that planets began to grow there. The new model predicts that the asteroid belt never contained much mass in bodies like the currently observed asteroids. “This presents the planetary science community with a testable prediction between this model and previous models that can be explored using data from meteorites, remote sensing, and spacecraft missions,” said Bottke. This work complements the recent study published in Nature by Levison, Kretke, and Martin Duncan (Queen’s University), which demonstrated that pebbles can form the cores of the giant planets and explain the structure of the outer solar system. Combined, the two works present the means to produce the entire solar system from a single, unifying process. “As far as I know, this is the first model to reproduce the structure of the solar system — Earth and Venus, a small Mars, a low-mass asteroid belt, two gas giants, two ice giants (Uranus and Neptune), and a pristine Kuiper Belt,” said Levison. The article, “Growing the Terrestrial Planets from the Gradual Accumulation of Sub-meter Sized Objects,” is published online by PNAS. Authors H.F. Levison, K.A. Kretke, K. Walsh, and W. Bottke are all of Southwest Research Institute’s Space Science and Engineering Division. This work was supported by the NASA Solar System Exploration Research Virtual Institute (SSERVI) through institute grant number NNA14AB03A.
The Senate of the United States shall be composed of two Senators from each State, elected by the people thereof, for six years; and each Senator shall have one vote. The electors in each State shall have the qualifications requisite for electors of the most numerous branch of the State legislatures. When vacancies happen in the representation of any State in the Senate, the executive authority of such State shall issue writs of election to fill such vacancies: Provided , That the legislature of any State may empower the executive thereof to make temporary appointments until the people fill the vacancies by election as the legislature may direct. This amendment shall not be so construed as to affect the election or term of any Senator chosen before it becomes valid as part of the Constitution. The Seventeenth Amendment is the only constitutional amendment to change the fundamental structure of the government as originally drafted in the Constitution. The Seventeenth Amendment increased the American public’s ability to control the federal government, because it granted voters the opportunity to directly elect their representatives to the Senate. Before the amendment was ratified in 1913, state legislatures chose senators. When the Constitution was written in 1787, many citizens wanted a “loose” union between the former colonies. This left the states with considerable powers to rule themselves as they wished. Under the original terms of the Constitution, the congressional houses divided the government’s power between the people and the states. Members popularly elected to the House of Representatives represented the American people, and states chose senators to represent them. The Seventeenth Amendment shifted the division of power in the government and gave the voters direct control over who represented their state. Submitted by Congress to the states on May 13, 1912. Ratified by the required three-fourths of states (thirty-six of forty-eight) by April 8, 1913, and by nine more states by March 9, 1922. Declared to be part of the Constitution on May 31, 1913. Massachusetts, May 22, 1912; Arizona, June 3, 1912; Minnesota, June 10, 1912; New York, January 15, 1913; Kansas, January 17, 1913; Oregon, January 23, 1913; North Carolina, January 25, 1913; California, January 28, 1913; Michigan, January 28, 1913; Iowa, January 30, 1913; Montana, January 30, 1913; Idaho, January 31, 1913; West Virginia, February 4, 1913; Colorado, February 5, 1913; Nevada, February 6, 1913; Texas, February 7, 1913; Washington, February 7, 1913; Wyoming, February 8, 1913; Arkansas, February 11, 1913; Maine, February 11, 1913; Illinois, February 13, 1913; North Dakota, February 14, 1913; Wisconsin, February 18, 1913; Indiana, February 19, 1913; New Hampshire, February 19, 1913; Vermont, February 19, 1913; South Dakota, February 19, 1913; Oklahoma, February 24, 1913; Ohio, February 25, 1913; Missouri, March 7, 1913; New Mexico, March 13, 1913; Nebraska, March 14, 1913; New Jersey, March 17, 1913; Tennessee, April 1, 1913; Pennsylvania, April 2, 1913; Connecticut, April 8, 1913. Origins of the Seventeenth Amendment When statesmen gathered at Independence Hall in Philadelphia in 1787, they intended to alter the Articles of Confederation and provide the framework for the new nation’s government. The Articles of Confederation outlined a country united by a weak federal government and strong states. However, the statesmen soon realized that a much stronger central government was needed in order to keep the union stable and at peace. They defined this new government in the Constitution. Checks and balances Having just won freedom from the unresponsive British monarchy, the American statesmen worked to create a responsive government of and by the people. Some statesmen feared that giving the public too much authority would subject the government to popular whims and lead to chaos and instability. These statesmen favored a government that included a system of checks and balances between the governing bodies. The checks and balances would protect the American people from themselves by allowing separate parts of the government to discuss an issue before committing the country to action. The resulting Constitution detailed such a system of checks and balances. In essays known as The Federalist, Alexander Hamilton, James Madison, and John Jay wrote that the checks and balances would protect the government and American people from “the blow meditated by the people against themselves, until reason, justice, and truth can regain their authority over the public mind.” Protecting the people from themselves. The Senate’s primary role in the new government was to provide a check on all legislation passed by Congress, or the lower house. It could also reject any treaty or political appointment initiated by the president. The framers of the Constitution reasoned that citizens were qualified to make good decisions about their representatives in the state legislatures. However, the framers felt those same citizens would not make good choices of senators to represent their state in the federal government. The framers feared that citizens might vote for undeserving politicians or that corrupt political interest groups who would take “advantage of the [indifference], the ignorance, and the hopes and fears of the unwary and interested.” The framers of the Constitution considered the Senate an anchor against such corruption. In addition, the framers considered a Senate chosen by state legislatures a good balance against a popularly elected Congress. The framers thought it unlikely that special interest groups would gain control if both houses were elected by different means. This would let the Senate and the Congress form a barrier against an interest group gaining control of legislative power. During the deliberations in Independence Hall, only James Wilson of Pennsylvania argued for the direct election (election by the people) of senators. But the decision to let state legislatures vote for senators was eventually adopted unanimously. The “Great Compromise.” The number of senators serving in the Senate was determined during the Federal Convention of 1787 in what is known as the “Great Compromise.” Individual states differed in size. They argued about how to select government representatives: depending on the state’s population or regardless of size. They compromised and created the House of Representatives. The House would represent the people of the states, and the Senate would represent the individual states. Two senators would represent each state. Large and small states alike were satisfied with the arrangement and never again competed for more representation at the federal level. Corruption in the Senate Until the end of the Civil War (1861-65), the U.S. Senate enjoyed a reputation as the “greatest deliberative (thinking) body in the Western world.” The Senate garnered the attention and praise of influential foreign dignitaries. After the Civil War, however, the Senate turned into a place where the interests of big businesses soon carried more weight than reasoned debate. After the Civil War, great wealth was generated for a small group of businessmen as industries combined and started to serve national markets. Railroads and other transportation companies were among those corporations that grew to serve a national market. Never before had so few companies controlled so much economic power. For example, the areas where railroads laid their tracks enjoyed more jobs, more tourists, and more access to markets than other areas. Big business takes over State legislatures quickly learned that they could use senatorial seats to gain favor with the influential businesses that brought wealth and jobs to their states. In turn, businessmen pursued senatorial seats when federal regulations started to impose limits on, or to provide benefits to, their businesses. In large cities such as Chicago, Kansas City, and New York City political bosses (businessmen who exerted a controlling force on political decisions) soon gained enough power to influence senatorial elections. Political bosses and other influential politicians and businessmen bought the votes of state legislatures or strong-armed (bullied) them, to effectively control who became a U.S. senator. If a candidate was not favored, wealthy businessmen or lobbyists for certain industries gave huge sums of money to all candidates. The money did not support a certain political viewpoint. Instead, the winning candidate was obligated to whoever gave him the money. Senator Chauncey Mitchell Depew was a prime example of this corruption. David Graham Phillips (see sidebar) wrote a muckraking article. He called Senator Depew “an ideal lieutenant for a plutocrat” and exposed Senator Depew’s connection with furthering the interests of the powerful William H. Vanderbilt family. (The Vanderbilts’ tried to push through legislation that would benefit their railroad between New York City and Buffalo, New York.) Phillips also linked corrupt senators and businessmen to legislation concerning beef inspection, food and drug purity standards, railroad regulations, and sugar subsidies (government granted money), among other things. Muckrakers expose big business The American public was not blind to the Senate’s corruption. Articles appeared in magazines and newspapers, telling tales of the corruption. The stories of the purchase of senatorial votes by political interest groups or wealthy businessmen were fantastic and sometimes exaggerated. However, they generated intense scrutiny of the Senate, and the voting records of its members. Skeptical people rejected the notion that many senators could be bribed, but they were also not blind to the number of millionaires in the Senate. The “Millionaire’s Club.” The Senate became known as “The Millionaire’s Club,” “The Rich Man’s Club,” and the “House of Dollars.” Among the millionaire senators were Philetus Sawyer of Wisconsin, who made his millions in lumber; California railroad magnate Leland Stanford; Arthur Gorman of Maryland, who ran the Chesapeake and Ohio Canal Company; and Nelson Aldrich of Rhode Island, who made millions in banking. The opinion that the Senate represented corporate wealth grew in popularity. State legislatures hinder senatorial elections. Corporate influence in the Senate was not the only problem with the legislative body. In the early 1800s, some states provided senators with written instructions for how they should vote on certain issues. Some senators resigned when their own beliefs did not coincide with their state’s wishes. Among them were John Quincy Adams in 1807 and John Tyler of Virginia in 1836. By the mid-1800s, states stopped making rules for how senators voted. The influence of political parties in senatorial elections. Some historians trace the influence of political parties in senatorial elections back to the first one in 1789. Early political parties represented different factions or commercial interests within the state. By the mid-1800s, political parties were more nationally organized and dominated senatorial elections. In some states, political parties would hold state conventions to choose their party’s nominee. The party’s representatives in the state legislature then promised to vote for their party’s nominee. Sometimes politicians were unwilling to vote for someone of another political party, and debates in state legislatures dragged on for months. Occasionally, this left states without representatives during some Congresses. George H. Haynes, a respected Senate historian, wrote that the sessions would sometimes degenerate into “riotous demonstrations more appropriate to a prize-fight than to a senatorial election.” In one particularly colorful instance, the Missouri state legislators threw fists, desks, and books at each other. The fight erupted over whether to break the wall clock, and allow deliberations over senatorial selections to continue after the scheduled hour of adjournment (closing). One member finally broke the clock by hurling ink bottles at it. “It is ridiculous to suggest that amid scenes like these the choice of a senator retains anything of the character of an exercise of cool judgment.” These types of sessions resulted in poor decisions, or no decisions at all. The most extreme case of state legislature paralysis was in Delaware. The Delaware state legislature’s inability to elect even one senator left the state without any representation in the federal government between 1901 and 1903. Other states succeeded in electing only one senator or elected senators after the Senate had already started a session. A Push for Direct Election of Senators Between 1826 and 1912, more than 197 resolutions for direct election of senators were introduced in the House of Representatives. Only six received the necessary two-thirds majority vote needed to reach the Senate. Once in the Senate, all of these resolutions were ignored. In 1910, the only resolution ever debated lost by a narrow margin. As the Senate continued to ignore the public’s will over the years, reformers of the election process were forced to unite across political parties. Though movement in the federal legislature was slow, political pressures to change the election process for senators gradually built momentum. As early as the 1880s and 1890s, reform advocates (supporters) declared that[MM1] “special interests had conspired to hold the Senate hostage,” and “the documentation they presented to the public painted a horrifying picture of a widespread network of corrupt bargains, in which wealth and power were exchanged for influence and votes.” Political parties made direct election of senators a part of their presidential platforms. It started with the People’s Party between 1892 and 1904 and was followed by the Democratic Party in 1900 and 1904 and then the Prohibition Party in 1904. The Pennsylvania legislature proposed a second Constitutional Convention in 1900. By 1905, thirty-one states had either passed referendums proposing that Congress consider a constitutional amendment or otherwise voiced their support for direct election. With the publication of muckraking articles such as David Graham Phillips’s “The Treason of the Senate” that appeared in Cosmopolitan Magazine in 1906, public interest in direct election mounted. Demand for reform was so great that the 57th Congress printed an additional 5,000 copies of the Senate committee report on direct election. A report of the Senate conducted during the 58th Congress (1903-1905) was published in 1906. It revealed that “One senator out of every three owes his election to his personal wealth, to his being the candidate satisfactory to what is coming to be called the ‘System,’ or to his expertness in political manipulation—qualifications which make their usefulness as members of the dominant branch of Congress decidedly open to question.” In 1906 alone, nine resolutions for direct election were put before Congress. The American public agreed that direct election would free the Senate from corrupting influences. In 1911, Indiana representative John A. Adair argued in Congress that direct election would fill the Senate with people with “rugged honesty, recognized ability, admitted capacity, and wide experience.” Popular election of Senators without an amendment Reformers looked for alternatives when they were unable to pass amendment proposals through the Senate or to gain enough support for a second Constitutional Convention. To appease popular pressure for direct election of senators, some states invented a new voting method to allow the public to select senators. The new method was called a primary election. Oregon developed the first primary elections between 1901 and 1904. With this method, people voted for candidates in primary elections. The state legislatures would then officially elect the primary’s winner to the Senate. Candidates pledged to uphold the popular primary elections during their tenure. Primary elections quickly delivered the desired reduction in deadlock and corruption in senatorial elections. An Oregon paper reported in 1907 that “On the first ballot, in twenty minutes, we elected two Senators, without boodle, or booze, or even a cigar!” By 1910, fourteen of the thirty states used primary elections to select senators. Muckraking and the Seventeenth Amendment Muckraking is a form of journalism. It was used in the early 1900s and boldly attempted to reveal some essential truth about public figures, political issues, or institutions. Muckrakers published articles that exposed corruptions in American government. According to George E. Mowry and Judson A. Grenier, muckrakers followed the advice of a biblical passage from St. John: “And ye shall know the truth and the truth shall make you free.” The articles enjoyed a great deal of popularity. Journalists such as Samuel Hopkins Adams, Ray Stannard Baker, Charles Edward Russell, Upton Sinclair, and Ida M. Tarbell became household names. Their vivid articles reached the readers of magazines such as McClure’s,Collier’s, the American, and Cosmopolitan. Muckrakers were primarily concerned with exposing the privileges that money bought in political life. The Industrial Revolution had created great wealth, but only for a very few. Muckrakers worked hard to show just how much influence this new wealth affected the government. The villains in the thousands of muckraking stories were greedy businessmen or what muckrakers called “predatory wealth.” To stop these corrupt businessmen, muckrakers called for greater democracy. They considered it “the inevitable sequence of widespread intelligence.” Publishers also made greater democracy a rallying cry. As early as February 5, 1899, William Randolph Hearst’s newspapers made the direct election of U.S. senators a firm editorial quest. The height of the Muckraking Era came in 1906 when David Graham Phillips published “The Treason of the Senate” articles in Cosmopolitan. In his articles, Phillips examined the corruption of the Senate. One of the most important issues Phillips tackled was the millionaires in the Senate. Charles Edward Russell was a fellow muckraker who believed the Senate was made up of senators who used their seats to make millions, Russell called the Senate a house made up of “butlers for industrialists and financiers.” Phillips reported that as many as twenty-five senators were millionaires at the time of his writing. (Some scholars estimate the number was actually closer to ten.) Powerful businessmen with great interest in creating laws to protect their companies often sought election. Phillips wondered whether it was better to elect these millionaires or to have them bribe senators. Phillips profiled the careers of twenty-one senators and declared that senatorial selection was based on private wealth and power in party organizations. Phillips called senators “grafters,” “bribers,” and “perjurers,” and exposed decisions in which senators favored the interests of their corporate backers over the interests of the public. President Roosevelt grew more frustrated with each monthly muckraking installment. He told the Post editor: “Phillips takes certain facts that are true in themselves, and by ignoring utterly a very much larger mass of facts that are just as true and just as important, and by downright perversion of truth both in the way of misstatement and of omission, succeeds in giving a totally false picture … [The articles] give no accurate guide for those who are really anxious to war against corruption, and they do excite a hysterical and ignorant feeling against everything existing, good or bad.” But Phillips responded to critics by stating that “these articles have been attacked, but their facts—the facts of the treason of the Senate, taken from the records—have not been attacked. Abuse is not refutation (a proof of being wrong); it is confession.” He added, “The exposed cry out that these exposures endanger the Republic. What a ludicrous inversion—the burglar shouting that the house is falling because he is being ejected from it? The Republic is not in danger; it is its enemies that are in danger.” In the end, “The Treason of the Senate” and other muckraking articles stirred what President Roosevelt labeled “a revolutionary feeling in the country.” The muckraking articles were read by hundreds of thousands of people and helped reformers generate popular consensus for the eventual ratification of the Seventeenth Amendment. Race and Ratification Another issue complicated the political power struggle over senatorial elections between the American public and corporate interests: the power of African American votes. Black votes could influence senatorial elections, and this was why some southern senators were reluctant about direct election proposals. The South had a long history of denying blacks the right to vote. The Fourteenth Amendment established equal rights for all citizens. The Fifteenth Amendment granted people of all races and colors the right to vote. Even after these amendments were ratified, southern states implemented poll taxes and literacy tests to keep as many blacks as possible from voting. To ensure that whites retained control at the polls, southern Democrats insisted that every proposal for direct election gave states the authority to regulate elections. The race rider A proposal for the first vote on direct election in the full Senate emerged from the Senate’s Judiciary Committee in 1910. The proposal had a provision that proved so controversial that it doomed the proposal from a passing vote. The provision was nicknamed the “race rider.” The race rider ensured that states would control election regulations. It had been added in the Judiciary Committee as a compromise to win Democratic votes. But when senators in the full Senate debated the proposal, many balked at the race rider. They declared it would effectively reverse the Fifteenth Amendment by allowing states to racially discriminate against some voters. Chauncey Depew of New York opposed it, saying that passing the proposal would be “deliberately voting to undo the results of the Civil War.” The direct election amendment failed to pass in the Senate in 1910, but several senators encouraged continued discussion of the issue. Over the next year, the House and the Senate continued to debate the race rider. The rider was the only real obstacle to the direct election proposal passing in both houses. After two months of intense debate, the Senate almost passed the proposal without the race rider on February 28, 1911. It was a fifty-four to thirty-three margin (four not voting)—just four votes short of the needed two-thirds majority. Congress went back and forth over the addition or removal of the race rider. The House voted on the amendment issue again in April 1911. The amendment passed with a race rider by a margin of 296 to 16 (70 not voting). The Senate discarded the race rider and passed the revised proposal with the required two-thirds supermajority by a margin of 64 to 24 (three not voting). Instead of quickly passing the revised version of the proposal, the House entered into yet further debate. Representative Walter Rucker proposed to abandon the race rider on May 13, 1912. The House voted again, and this time passed the revised amendment proposal by a margin of 238 to 39 (five voted “present,” and 110 not voting). Unlike the congressional houses, the states did not wrangle over the direct-election amendment. The amendment swiftly passed through state legislatures in less than a year and became the Seventeenth Amendment on April 8, 1913. Although the amendment decreased the power of the states, the state legislatures seemed more than happy to pass it. With the use of primary elections, many of the U.S. senators had already been selected by popular vote. To many, the ratification of the Seventeenth Amendment simply formalized a practice of direct election that was already widely used. Although the Seventeenth Amendment was ratified without a race rider, many states found ways to limit black voters from participating in senatorial elections. To keep minorities from the polls, southern states instituted poll taxes and literacy tests. These taxes made voting too expensive for poorer blacks. The literacy tests tested people’s knowledge of the meaning of the state constitution. By the 1920s, the “white primary” was the main way southern states blocked blacks and other minorities from voting for senators. The white primary was named as such after southern political parties made exclusive rules to forbid minority membership. In the South, the Democratic Party did not allow black members. To make matters even more unfair, state legislatures only allowed party members to vote in primary elections. Up until the 1960s, the Democratic Party was so strong in the South that Democratic nominees nearly always won election. Sometimes the Republican Party did not even bother to select a party candidate. Because the Democratic Party dominated southern politics, whites were the only people allowed to vote in most primary elections. It took several decades and several court cases to stop discriminatory voting practices in the South. One of the most influential cases was Chapman v. King in 1946. In Chapman, the Fifth Circuit Court of Appeals determined that the Democratic primary in Georgia violated black people’s constitutional right to vote. The court ruled that since the Georgia state legislature accepted the Democratic primary candidate, the state also endorsed the Democratic Party’s discriminatory rules against black participation. By doing this, the Georgia state legislature denied African Americans their constitutional right to vote for their own representatives. Though the court’s decision was favorable, southern blacks endured many more decades of discrimination before they could fully exercise their right to vote for U.S. senators. The white primary effectively stopped black voters from voting until the Voting Rights Act was passed in 1965. The Voting Rights Act came at the crest of the civil rights movement in the 1960s. The Twenty-fourth Amendment abolished the poll tax in 1962, but it did not end all obstacles for minority voters. The amendment helped rally support for further guarantees for minorities. President Lyndon B. Johnson supported both the Civil Rights Act of 1964 and the Voting Rights Act. The Civil Rights Act made discriminatory practices in employment, education, and public places illegal The Voting Rights Act similarly rendered illegal the remaining deterrents that stopped minorities from voting. The Progressive Amendments The Seventeenth Amendment is tied to a time between when Americans united and gave themselves more control over their public and private lives. The time is called the Progressive Era, and it lasted from about 1900 to 1920. The “Progressive” amendments are: the Sixteenth Amendment, the Seventeenth Amendment, and the Nineteenth Amendment. (The Sixteenth Amendment established a federal income tax, and the Nineteenth Amendment gave women the right to vote.) The Seventeenth Amendment made the federal government more democratic. It was the greatest alteration to the workings of the state and federal governments since the Civil War amendments—the Thirteenth, Fourteenth, and Fifteenth. (Together, they abolished slavery, established equal rights for all citizens, and granted all races the right to vote.) New ways to use amendments The passage of the Sixteenth and the Seventeenth Amendments in 1913 ushered in a wave of new thinking about the purpose of constitutional amendments. During the forty years following the Civil War, no constitutional amendments were ratified. Hundreds of amendments were proposed during those forty years. However, none could gain the two-thirds majority vote needed to ratify them. Politicians and social advocates questioned the usefulness of the Constitution, thinking it was inflexible. But when two constitutional amendments were ratified in the same year, constitutional amendments seemed like realistic solutions to a variety of social and political problems. The American public recognized the power constitutional amendments had to redirect the activities of government. Perhaps more important, Americans also realized their own ability to effectively change the Constitution to create a government that was responsive to current needs and values. The ratification of the Sixteenth and Seventeenth Amendments in the same year was “a political reaction to the great concentration of wealth and its alleged corrupting influence on the political system.” In reaction to the alterations in the country’s economic structure brought about by the Civil War, the American public demanded a more equal tax structure. The Sixteenth Amendment provided that tax structure and laid the groundwork for weakening the concentration of wealth in the country. But the public grew frustrated as it heard about the Senate debates over the Sixteenth Amendment. Senators who owned or were influenced by the large corporations tried to block the amendment. This situation drew the public’s attention to the great number of millionaires in the Senate. The public was angry that the wealthy senators would not pass a more equitable tax system, and they wondered how to make the senators more responsive to their views. Public outcry to the senators’ attempts to block the Sixteenth Amendment helped to push the Seventeenth Amendment through the Senate. The Sixteenth and Seventeenth Amendment paved the way for Progressive advocates to succeed in passing the other Progressive amendments. The other amendments were designed to take control of government from large corporations and give it back to the people. Constitutional amendments fixed more than just governmental procedural problems. They were ways to restrict alcohol use, grant voting rights for women, and regulate child labor. Effects of the Seventeenth Amendment The Seventeenth Amendment succeeded in some ways, but failed in others. The Seventeenth Amendment made the Senate more responsive to the American public. The American public gained greater access to senators and influence in senators’ decisions on many issues of national and international importance. Senators, who must run for election and reelection, discovered that the best policy is to aim for openness and responsiveness to public opinion. One area in which the Senate has been accused of being too responsive to the public is in its handling of presidential appointments to the Supreme Court. Public opinion heavily influenced Senate reactions during two appointment processes. First, in 1987 when President Ronald Reagan (1911–2004; served 1981–89) nominated Robert H. Bork, and next in 1991 when President George H. Bush (1924–; served 1989–93) successfully appointed Clarence Thomas to the Supreme Court. The Senate’s Judiciary Committee considers Supreme Court candidates and can also generate strong public reactions. News stories triggered heated public opposition to the Bork and Thomas nominations. The network news reported about Judge Bork’s perceived (by Democrats) controversial legal opinions and theories on constitutional law, and there were televised hearings of sexual harassment charges brought against Clarence Thomas by law professor Anita Hill, a former employee who worked for Thomas at the Equal Employment Opportunity Commission (EEOC). Although many anticipated that direct election would make U.S. senators prone to popular whims, the Seventeenth Amendment also changed the Senate in unexpected ways. Rather than ridding the Senate of millionaires, the amendment heightened the importance of money in the Senate. Compared to the ten millionaires in the Senate in 1900, there were more than twenty-five by the mid-1990s. Senators must constantly campaign in order to raise the funds needed for reelection. It is estimated that senators must raise approximately $15,000 for each week of their six-year terms. Candidates in California have spent more than $10 million to secure a Senate seat. By the 1990s, an average expenditure per seat in other states had reached $5 million. Criticism in the early 2000s of the Seventeenth Amendment In the early 2000s, many groups and commentators have called for the repeal of the Seventeenth Amendment, believing that modern elections have become too costly and allowed special interest groups to exert too much influence over the process. In the view of these critics, the Seventeenth Amendment reduced the influence of states and expanded the role of special-interest groups. Others believe that the Seventeenth Amendment has elevated the power of the federal government too far above that of the states. The reasoning is that before the Seventeenth Amendment, state legislatures had the power to decide who served in the Senate. Such senators naturally would look more to state interests. The Seventeenth Amendment removed the direct pipeline from state legislatures to the U.S. Senate. Some have argued that another impact of the Seventeenth Amendment: It causes federal courts to be more likely to strike down state laws. Professor Donald J. Kochan wrote in 2003: “there is substantial empirical evidence that suggests that the Seventeenth Amendment may have altered the relationship between state legislatures and federal courts.” In April 2004, Senator Zell Miller (Democrat, Georgia) introduced Senate Joint Resolution 135, which called for the repeal of the Seventeenth Amendment. In introducing the measure to his colleagues on the Senate floor, Miller said: Perhaps, then, the answer is a return to the original thinking of those wisest of all men, and how they intended for this government to function. Federalism, for all practical purposes, has become to this generation of leaders, some vague philosophy of the past that is dead, dead, dead. It isn’t even on life support. The line on that monitor went flat some time ago. You see, the reformers of the early 1900s killed it dead and cremated the body when they allowed for the direct election of U.S. senators. Up until then, senators were chosen by State legislatures, as James Madison and Alexander Hamilton had so carefully crafted. Direct elections of senators, as great and as good as that sounds, allowed Washington’s special interests to call the shots, whether it is filling judicial vacancies, passing laws, or issuing regulations. The State governments aided in their own collective suicide by going along with that popular fad at the time. … The Seventeenth Amendment was the death of the careful balance between State and Federal Government. As designed by that brilliant and very practical group of Founding Fathers, the two governments would be in competition with each other and neither could abuse or threaten the other. The election of senators by the state legislatures was the lynchpin that guaranteed the interests of the states would be protected. Efforts have also been made in state legislatures, including Montana, to repeal the Seventeenth Amendment. FOR MORE INFORMATION Chopra, Pram. Supreme Court versus the Constitution: A Challenge to Federalism. New Delhi: Sage Publications India, 2006. Haskell, John. Direct Democracy or Representative Government?: Dispelling the Populist Myth. Boulder, CO: Westview Press, 2001. Holzer, Henry Mark. Supreme Court Opinions of Clarence Thomas (1991–2006): A Conservative’s Perspective. Jefferson, NC: McFarland, 2007. Merida, Kevin, and Michael Fletcher. Supreme Discomfort: The Divided Soul of Clarence Thomas. New York: Doubleday, 2007. Rossum, Ralph A. Federalism, the Supreme Court, and the Seventeenth Amendment: The Irony of Constitutional Democracy. Lexington, MA: Lexington Press, 2001. Kochan, Donald J. “State Laws and the Independent Judiciary: An Analysis of the Effects of the Seventeenth Amendment on the Number of Supreme Court Cases Holding State Laws Unconstitutional.” Albany Law Review, vol. 66, (2003). Rossum, Ralph A. “The Irony of Constitutional Democracy: Federalism, the Supreme Court, and the Seventeenth Amendment.” San Diego Law Review, vol. 36, (1999). The Center for Constitutional Studies. (accessed August 22, 2007) . Dean, John W. "The Seventeenth Amendment: Should It Be Repealed?" Findlaw’s Writ , (accessed August 22, 2007). “Repeal the Seventeen Amendment.” (accessed August 22, 2007). “Things that Are Not in the U.S. Constitution.” The U.S. Constitution Online. http:/www.usconstitution.net/constnot.html (accessed August 22, 2007).
Presentation on theme: "Week 1: Fundamentals of Sound Sound Levels and the Decibel"— Presentation transcript: 1 Week 1: Fundamentals of Sound Sound Levels and the Decibel AcousticsWeek 1:Fundamentals of SoundSound Levels and the Decibel 2 Fundamentals of SoundSound can be defined as a wave motion in air or other elastic media (stimulus) or as that excitation of the hearing mechanism that results in the perception of sound (sensation).Frequency is a characteristic of periodic waves measured in hertz (cycles per second), readable on an oscilloscope or frequency counter. The ear perceives a different pitch for a quiet tone than a loud one. The pitch of a low-frequency tone goes down, while the pitch of a high-frequency tone goes up as intensity increases. We cannot equate frequency and pitch, but they are analogous.The same situation exists between intensity and loudness. The relationship between the two is not linear.Similarly, the relationship between waveform (or spectrum) and perceived quality (or timbre) is complicated by the functioning of the hearing mechanism. 3 Sine WavesThe sine wave is a basic waveform closely related to simple harmonic motion. Vibration or oscillation is possible because of the elasticity of the spring and the inertia of the weight. Elasticity and inertia are two things all media must possess to be capable of conducting sound. 4 Sine Wave LanguageThe easiest value to read is the peak-to-peak value (of voltage, current, sound pressure, etc.). If the wave is symmetrical, the peak-to-peak value is twice the peak value. Another common way to measure the sine wave is using RMS (Root Mean Square) values. The other two scales of measurement used are peak and average. The picture to the right shows four different ways that a sine wave’s amplitude can be measured and the equations used for conversions. 5 Propagation of SoundIf an air particle is displaced from its original position, elastic forces of the air tend to restore it to its original position. Because of the inertia of the particle, it overshoots the resting position.Sound is readily conducted in gases, liquids, and solids such as air, water, steel, etc., which are all elastic media.Without a medium, sound cannot be propagated. Outer space is an almost perfect vacuum; no sound can be conducted due to the absence of air.Particles of air propagating a sound wave do not move far from their undisplaced positions. 6 Forms of Particle Motion There is more than a million molecules in a cubic inch of air.The molecules crowded together represent areas of compression and the sparse areas represent rarefactions. 7 Sound in Free SpaceThe intensity of sound decreases as the distance to the source is increased.Doubling the distance reduces the intensity to 1/4 the initial value, tripling the distance yields 1/9, increasing the distance four times yields 1/16 the initial intensity.The inverse square law states that the intensity of sound in a free field is inversely proportional to the square of the distance from the source. 8 Wavelength, Period, & Frequency The wavelength is the distance a wave travels in the time it takes to complete one cycle. The period is the time it takes a wave to complete one cycle. The frequency is the number of cycles per second (Hertz). 9 Wavelength and Frequency Formulas for calculating wavelength and frequency. The speed of sound in air is about 1,130 feet per second at normal temperature and pressure.Two graphical approaches for an easy solution to the above equations. 10 Complex WavesThe sine wave with the lowest frequency (f1) is called the fundamental, the one with twice the frequency (f2) is called the second harmonic, and the one three times the frequency (f3) is the third harmonic. Harmonics are whole number multiples of the fundamental frequency. 11 PhasePhase is the time relationship between waveforms. Each waveform is lagging the previous by 90 degrees. 12 Combinations of Waveforms Combinations of waveforms that are not in phase. The difference in waveshapes is due entirely to the shifting of the phase of the harmonics with respect to the fundamental. 14 SpectrumThe audio or frequency spectrum of the human ear is about 20 Hz to 20 kHz. The spectrum tells how the energy of the signal is distributed in frequency. For the ideal sine wave, all the energy is concentrated at one frequency. All other types of waveforms have more than one frequency present.The diagram shows various types of waveforms and their harmonic content. The sine, triangle, and square waves are known as period waves due to their cyclic pattern. 15 Sound Levels and the Decibel Levels in decibels make it easy to handle the extremely wide range of sensitivity in human hearing.The threshold of hearing matches the ultimate lower limit of perceptible sound in air.A level in decibels is a convenient way of handling the billion fold range of sound pressures to which the ear is sensitive without getting bogged down in a long string of zeros. 16 Ratios vs. DifferenceRatios of pressure seem to describe loudness changes better than difference in pressure.Ernst Weber (1834), Gustaf Fechner (1860), Hermann von Helmholtz (1873), and other researchers pointed out the importance of ratios.Ratios of stimuli come closer to matching up with human perception than do differences of stimuli.Ratios of powers, intensities, sound pressure, voltage, or anything else are unitless. This is important because logarithms can be taken only of unitless numbers. 17 Here are three different ways that numbers can be expressed Handling NumbersHere are three different ways that numbers can be expressed 46 DecibelsA level is a logarithm of a ratio of two like-power quantities.A level in decibels is ten times the logarithm to the base 10 of the ratio of two power quantities. 47 Decibels Sound pressure is proportional to (sound power)2. The squaring of the sound power results in the equation SPL = 20 log (p1/p2) instead of 10 log. 48 Sound Pressure Level (SPL) Sound pressure is usually the most accessible parameter to measure in acoustics, as voltage is for electronic circuits. For this reason, the Equation (2-3) form is more often encountered in day-to-day technical work. 49 Reference LevelsA sound level meter is used to read sound pressure levels. A sound level meter reading is a certain sound pressure level, 20 log (p1/p2). For sound in air, the standard reference pressure is 20 μPa (micropascal). The μPa is a very minute sound pressure and corresponds closely to the threshold of human hearing.The equations to the right show how to convert SPL (in dB) to μPa. 50 Log-to-Exponent Conversion Here’s another quick look at the math used to convert logs to exponents: 66 Measuring Sound Pressure Level Sound level meters usually offer a selection of weighting networks designated A, B, and C. Network selection is based on the general level of sounds to be measured. For SPLs of 20 – 55 dB, use network A. For 55 – 85 dB use network B, and for 85 – 140 dB, use network C. Your consent to our cookies if you continue to use this website.
Surface Tension Theory_E thinkiit lecture notes iit physics... SURFACE TENSION EXPLANATION OF SOME OBSERVED PHENOMENA 1. Lead balls are spherical in shape. Rain drops and a globule of mercury placed on glass plate are spherical. Hair of a shaving brush/painting brush, when dipped in water spread out, but as soon as it is taken out. Its hair stick together. A greased needle placed gently on the free surface of water in a beaker does not sink. Similarly, insects can walk on the free surface of water without drowning. Bits of Camphor gum move irregularly when placed on water surface. SURFACE TENSION Surface Tension is a property of liquid at rest by virtue of which a liquid surface gets contracted to a minimum area and behaves like a stretched membrane. Surface Tension of a liquid is measured by force per unit length on either side of any imaginary line drawn tangentially over the liquid surface, force being normal to the imaginary line as shown in the figure. i.e. Surface tension (T) = Total force on either of the imginary line (F) Length of the line () Units of Surface Tension. In C.G.S. system the unit of surface tension is dyne/cm (dyne cm -1) and Sl system its units is Nm -1 A ring is cut form a platinum tube of 8.5 cm internal and 8.7 cm external diameter. It is supported horizontally from a pan of a balance so that it comes in contact with the water in a glass vessel. What is the surface tension of water if an extra 3.97 g weight is required to pull it away from water? (g = 980 cm/s2). T T Cross section The ring is in contact with water along its inner and outer circumference; so when pulled out the total force on it due to surface tension will be F = T (2 r1 + 2 r2) So, mg T = 2(r r ) 1 2 [ F = mg] 3.97 980 = 72.13 dyne/cm 3.14 (8.5 8.7) EXCESS PRESSURE INSIDE A LIQUID DROP AND A BUBBLE 1. Inside a bubble : Consider a soap bubble of radius r. Let p be the pressure inside the bubble and pa outside. The excess pressure = p – pa. Imagine the bubble broken into two halves, and consider one half of it as shown in Fig. Since there are two surfaces, inner and outer, so the force due to surface tension is F = surface tension x length = T x 2 (circumference of the bubble) = T x 2 (2 T r) ... (1) The excess pressure (p - pa) acts on a cross-sectional area r2, so the force due to excess PHYSICS pressure is F = (p – pa) r2 .......... (2) The surface tension force given by equation (1) must balance the force due to excess pressure given by equation (2) to maintain the equilibrium. i.e. (p – pa) r2 = T × 2 (2 r) 4T = pexcess r above expression can also be obtained by equation of excess pressure of curve surface by putting R1 = R2. or (p – pa) = Inside the drop : In a drop, there is only one surface and hence excess pressure can be written as 2T = pexcess r 2T (p – pa) = = pexcess r (p – pa) = Inside air bubble in a liquid : A charged bubble : If bubble is charged, it's radius increases. Bubble has pressure excess due to charge too. Initially pressure inside the bubble 4T = pa + r 1 4T 2 For charge bubble, pressure inside = pa + r – , where surface is surface charge density.. 2 0 2 Taking temperature remains constant, then from Boyle's law 4T 2 4 4T 4 3 p 3 pa a r2 2 0 3 r2 r1 3 r1 = From above expression the radius of charged drop may be calculated. It can conclude that radius of charged bubble increases, i.e. r 2 > r1 Ex. 2 A minute spherical air bubble is rising slowly through a column of mercury contained in a deep jar. If the radius of the bubble at a depth of 100 cm is 0.1 mm, calculate its depth where its radius is 0.126 mm, given that the surface tension of mercury is 567 dyne/cm. Assume that the atmospheric pressure is 76 cm of mercury. The total pressure inside the bubble at depth h1 is (P is atmospheric pressure) = (P + h1 g) + 2T = P1 r1 and the total pressure inside the bubble at depth h2 is = (P + h2 g) + 2T = P2 r2 Now, according to Boyle’s Law P1V1 = P2V2 4 r13 , 3 Hence we get 2T 4 2T 4 3 3 (P h1 g) (P h2 g) r = 1 r1 3 r2 3 r2 2T 2T (P h1 g) r13 = (P h 2 g) 3 r1 r2 r2 4 r23 3 Given that : h1 = 100 cm, r1 = 0.1 mm = 0.01 cm, r2 = 0.126 mm = 0.0126 cm, T = 567 dyne/cm, P = 76 cm of mercury. Substituting all the values, we get h2 = 9.48 cm. PHYSICS THE FORCE OF COHESION The force of attraction between the molecules of the same substance is called cohesion. In case of solids, the force of cohesion is very large and due to this solids have definite shape and size. On the other hand, the force of cohesion in case of liquids is weaker than that of solids. Hence liquids do not have definite shape but have definite volume. The force of cohesion is negligible in case of gases. Because of this fact, gases have neither fixed shape nor volume. Examples. (i) Two drops of a liquid coalesce into one when brought in mutual contact because of the cohesive force. (ii) It is difficult to separate two sticky plates of glass wetted with water because a large force has to be applied against the cohesive force between the molecules of water. It is very difficult to break a drop of mercury into small droplets because of large cohesive force between mercury molecules. FORCE OF ADHESION The force of attraction between molecules of different substances is called adhesion. Examples. (i) Adhesive force enables us to write on the black board with a chalk. (ii) Adhesive force helps us to write on the paper with ink. (iii) Large force of adhesion between cement and bricks helps us in construction work. (iv) Due to force of adhesive, water wets the glass plate. (v) Fevicol and gum are used in gluing two surfaces together because of adhesive force. ANGLE OF CONTACT The angle which the tangent to the liquid surface at the point of contact makes with the solid surface inside the liquid is called angle of contact. Those liquids which wet the walls of the container (say in case of water and glass) have meniscus concave upwards and their value of angle of contact is less than 90° (also called acute angle). However, those liquids which don't wet the walls of the container (say in case of mercury and glass) have meniscus convex upwards and their value of angle of contact is greater than 90° (also called obtuse angle). The angle of contact of mercury with glass about 140°, whereas the angle of contact of water with glass is about 8°. But, for pure water, the angle of contact with glass is taken as 0°. SHAPE OF LIQUID MENISCUS When a capillary tube or a tube is dipped in a liquid, the liquid surface becomes curved near the point of contact. This curved surface is due to the two forces i.e. (i) due to the force of cohesion and (ii) due to the force of adhesion. The curved surface of the liquid is called meniscus of the liquid. Various forces acting on molecule A are: Force F1 due to force of adhesion which acts outwards at right angle to the wall of the tube. This force is represented by AB. Force F2 due to force of cohesion which acts at an angle of 45° to the vertical. This force is represented by AD. The weight of the molecule A which acts vertically downward along the wall of the tube. Since the weight of the molecule is negligible as compared to F 1 and F 2 and hence can be neglected. Thus, there are only two forces (F 1 and F2) acting on the liquid molecules. These forces are inclined at an angle of 135°. PHYSICS The resultant force represented by AC will depend upon the values of F 1 and F2. Let the resultant force makes an angle with F1. According to parallelogram law of vectors F2 sin135 º F2 / 2 tan = F F cos 135 º = = F1 F2 / 2 1 2 F2 2 F1 F2 Special cases : (i) If F2 = 2 F1, then tan = 90° Then the resultant force will act vertically downward and hence the meniscus will be plane or horizontal shown in figure (a). Example; pure water contained in silver capillary tube. If F2 < 2 F1, then tan is positive is acute angle Thus, the resultant will be directed outside the liquid and hence the meniscus will be concave upward shown in figure (b). This is possible in case of liquids which wet the walls of the capillary tube. Example ; water in glass capillary tube. If F2 > 2 F1, then tan is negative is obtuse angle. Thus, the resultant will be directed inside the liquid and hence the meniscus will be convex upward shown in figure (c). This is possible in case of liquids which do not wet the walls of the capillary tube. Example ; mercury in glass capillary tube. RELATION BETWEEN SURFACE TENSION, RADII OF CURVATURE AND EXCESS PRESSURE ON A CURVED SURFACE. Let us consider a small element ABCD (fig.) of a curved liquid surface which is convex on the upper side. R1 and R2 are the maximum and minimum radii of curvature respectively, They are called the 'principal radii of curvature' of the surface. Let p be the excess pressure on the concave side. 1 1 then p = T R R . If instead of a liquid surface, 1 2 we have a liquid film, the above expression will be 1 1 p = 2T R R , because a film has two surface. 1 2 EXCESS OF PRESSURE INSIDE A CURVED SURFACE 1. Plane Surface : If the surface of the liquid is plane [as shown in Fig.(a)], the molecule on the liquid surface is attracted equally in all directions. The resultant force due to surface tension is zero. The pressure, therefore, on the liquid surface is normal. Concave Surface : If the surface is concave upwards [as shown in Fig.(b)], there will be upward resultant force due to surface tension acting on the molecule. Since the molecule on the surface is in equilibrium, there must be an excess of pressure on the concave side in the downward direction to balance the resultant force of surface tension pA – pB = 2T . r Convex Surface : If the surface is convex [as shown in Fig.(c)], the resultant force due to surface tension acts in the downward direction. Since the molecule on the surface are in equilibrium, there must be an excess of pressure on the concave side of the surface acting in the upward direction to balance the downward resultant force of surface tension, Hence there is always an excess of pressure on concave side of a curved surface over that on the convex side. 2T pB – pA = r PHYSICS Ex. 3 A barometer contains two uniform capillaries of radii 1.44 × 10–3 m and 7.2 × 10–4 m. If the height of the liquid in the narrow tube is 0.2 m more than that in the wide tube, calculate the true pressure difference. Density of liquid = 103 kg/m3, surface tension = 72 × 10–3 N/m and g = 9.8 m/s2. Let the pressure in the wide and narrow capillaries of radii r1 and r2 respectively be P1 and P2. Then pressure just below the mensiscus in the wide and narrow tubes respectively are 2T 2T P1 and P2 r r2 1 [excess pressure = 2T ]. r 2T 2T = P1 r – P2 r = hg 1 2 True pressure difference = P1 – P2 Difference in these pressures 1 1 = hg + 2T r r 1 2 1 1 = 0.2 × 103 × 9.8 + 2 × 72 × 10–3 3 4 7.2 10 1.44 10 3 2 = 1.86 × 10 = 1860 N/m CAPILLARITY A glass tube of very fine bore throughout the length of the tube is called capillary tube. If the capillary tube is dipped in water, the water wets the inner side of the tube and rises in it [shown in figure (a)]. If the same capillary tube is dipped in the mercury, then the mercury is depressed [shown in figure (b)]. The phenomenon of rise or fall of liquids in a capillary tube is called capillarity. PRACTICAL APPLICATIONS OF CAPILLARITY 1. 2. 3. 4. 5. 6. 7. 8. The oil in a lamp rises in the wick by capillary action. The tip of nib of a pen is split up, to make a narrow capillary so that the ink rises upto the tin or nib continuously. Sap and water rise upto the top of the leaves of the tree by capillary action. If one end of the towel dips into a bucket of water and the Other end hangs over the bucket the towel soon becomes wet throughout due to capillary action. Ink is absorbed by the blotter due to capillary action. Sandy soil is more dry than clay. It is because the capillaries between sand particles are not so fine as to draw the water up by capillaries. The moisture rises in the capillaries of soil to the surface, where it evaporates. To preserve the moisture m the soil, capillaries must be. broken up. This is done by ploughing and leveling the fields Bricks are porous and behave like capillaries. CAPILLARY RISE (HEIGHT OF A LIQUID IN A CAPILLARY TUBE) ASCENT FORMULA Consider the liquid which wets the walls of the tube, forms a concave meniscus shown in figure. Consider a capillary tube of radius r dipped in a liquid of surface tension T and density p. Let h be the height through which the liquid rises in the tube. Let p be the pressure on the concave side of the meniscus and pa be the pressure on the convex side of the meniscus. The excess pressure 2T R Where R is the radius of the meniscus. Due to this excess pressure, the liquid will rise in the capillary tube till it becomes equal to the hydrostatic pressure hpg. Thus in equilibrium state. (p – pa) is given by (p – pa) = PHYSICS 2T = hpg R Let be the angle of contact and r be the radius of the capillary tube shown in the fig. Excess pressure = Hydrostatic pressure r OC 2T cos = cos or R = h = rg OA cos This expression is called Ascent formula. Discussion. (i) For liquids which wet the glass tube or capillary tube, angle of contact < 90°. Hence cos = positive. h = positive. It means that these liquids rise in the capillary tube. Hence, the liquids which wet capillary tube rise in the capillary tube. For example, water, milk, kerosene oil, patrol etc. A liquid of specific gravity 1.5 is observed to rise 3.0 cm in a capillary tube of diameter 0.50 mm and the liquid wets the surface of the tube. Calculate the excess pressure inside a spherical bubble of 1.0 cm diameter blown from the same liquid. Angle of contact = 0º. The surface tension of the liquid is T= (0.025 cm) (3.0 cm) (1.5 gm / cm3 ) (980 cm / sec 2 ) 2 = 55 dyne/cm. Hence excess pressure inside a spherical bubble = 4 55 dyne / cm 4T = = 440 dyne/cm2 . (0.5 cm) R For liquids which do not wet the glass tube or capillary tube, angle of contact > 90°. Hence cos = negative h = negative. Hence, the liquids which do not wet capillary tube are depressed in the capillary tube. For example, mercury. T, , and g are constant and hence h If two parallel plates with the spacing 'd' are placed in water reservoir, then height of rise 2T = hdg 1 . Thus, the liquid r rises more in a narrow tube and less in a wider tube. This is called Jurin's Law. 2T h = dg (v) If two concentric tubes of radius 'r 1' and 'r2' (inner one is solid) are placed in water reservoir, then height of rise T [2r1 + 2r2] = [r22 h – r12h] g 2T h = (r r )g 2 1 If weight of the liquid in the meniscus is to be consider : T cos × 2r = [r2h + 1 2 r × r] g 3 2T cos r h 3 = rg When capillary tube (radius, 'r') is in vertical position, the upper meniscus is concave and pressure due to surface tension is directed vertically upward and is given by p1 = 2T/R1 where R1 = radius of curvature of upper meniscus. The hydrostatic pressure p2 = h g is always directed downwards. If p1 > p2 i.e. resulting pressure is directed upward. For equilibrium, the pressure due to lower meniscus should be downward. This makes lower meniscus concave downward (fig.a). The 2T radius of lower meniscus R2 can be given by R = (p1 – p2). 2 If p1 < p2, i.e. resulting pressure is directed downward for equilibrium, the pressure due to lower meniscus should be upward. This makes lower meniscus convex upward (fig. b). 2T The radius of lower meniscus can be given by R = p2 – p1. 2 2T If p1 = p2, then is no resulting pressure. then, p1 – p2 = R = 0 or, R2 = i. e. lower surface will 2 be FLAT. (fig.c). (viii) Liquid between two Plates - When a small drop of water is placed between two glass plates put face to face, it forms a thin film which is concave outward along its boundary. Let 'R' and 'r' be the radii of curvature of the enclosed film in two perpendicular directions. Hence the pressure inside the film is less than the atmospheric pressure outside it by an T 1 1 amount p given by p = T and we have. p = . r r If d be the distance between the two plates and the angle of contact for water and glass, then, 1 d 2 from the figure, cos = or r 1 2 cos = . r d 1 2T in , we get p = cos . r d can be taken zero for water and glass, i.e. cos = 1. Thus the upper plate is pressed Substituting for downward by the atmospheric pressure minus acting on the upper plate is 2T . Hence the resultant downward pressure d 2T . If A be the area of the plate wetted by the film, the resultant d 2TA . d For very nearly plane surface, d will be very small and hence the pressing force F very large. Therefore it will be difficult to separate the two plates normally. force F pressing the upper plate downward is given by F = resultant pressure × area = A drop of water volume 0.05 cm 3 is pressed between two glass-plates, as a consequence of which, it spreads and occupies an area of 40 cm 2. If the surface tension of water is 70 dyne/cm, find the normal force required to seperate out the two glass plates in newton. 1 1 Pressure inside the film is less than outside by an amount, P = T r r , where r1 and r2 are the radii 1 2 of curvature of the meniscus. Here r 1 = t/2 and r2 = , then the force required to separate the two glassplates, between which a liquid film is enclosed (figure) is, F = P × A = 2AT , where t is the thickness t of the film, A = area of film. 2 ( 40 10 4 )2 (70 10 3 ) 2A 2 T 2A 2 T = At V 0.05 10 6 = 45 N A glass plate of length 10 cm, breadth 1.54 cm and thickness 0.20 cm weighs 8.2 gm in air. It is held vertically with the long side horizontal and the lower half under water. Find the apparent weight of the plate. Surface tension of water = 73 dyne per cm, g = 980 cm/sec2. Volume of the portion of the plate immersed in water is 1 10 × (1.54) × 0.2 = 1.54 cm3. 2 Therefore, if the density of water is taken as 1, then upthrust = wt. of the water displaced = 1.54 × 1 × 980 = 1509.2 dynes. Now, the total length of the plate in contact with the water surface is 2(10 + 0.2) = 20.4 cm, downward pull upon the plate due to surface tension = 20.4 × 73 = 1489.2 dynes resultant upthrust = 1509.2 – 1489.2 20 = 20.0 dynes = 980 = = 0.0204 gm-wt. apparent weight of the plate in water = weight of the plate in air – resultant upthrust = 8.2 – 0.0204 = 8.1796 gm Ans. A glass tube of circular cross-section is closed at one end. This end is weighted and the tube floats vertically in water, heavy end down. How far below the water surface is the end of the tube? Given : Outer radius of the tube 0.14 cm, mass of weighted tube 0.2 gm, surface tension of water 73 dyne/cm and g = 980cm/sec2. Let be the length of the tube inside water. The forces acting on the tube are : (i) Upthrust of water acting upward 22 = r2 × 1 × 980 = × (0.14)2 × 980 = 60.368 dyne. 7 (ii) Weight of the system acting downward = mg = 0.2 × 980 = 196 dyne. (iii) Force of surface tension acting downward = 2rT 22 × 0.14 × 73 = 64.24 dyne. 7 Since the tube is in equilibrium, the upward force is balanced by the downward forces. That is, 60.368 = 196 + 64.24 = 260.24. =2× = 4.31 cm. A glass U-tube is such that the diameter of one limb is 3.0 mm and that of the other is 6.00 mm. The tube is inverted vertically with the open ends below the surface of water in a beaker. What is the difference between the heights to which water rises in the two limbs? Surface tension of water is 0.07 nm–1. Assume that the angle of contact between water and glass is 0º. Suppose pressures at the points A, B, C and D be PA, PB, PC and PD respectively. The pressure on the concave side of the liquid surface is greater than that on the other side by 2T/R. An angle of contact is given to be 0º, hence R cos 0º = r or R = r PA = PB + 2T/r1 and PC = PD + 2T/r2 where r1 and r2 are the radii of the two limbs A But PA = PC B 2T 2T PB + = PD + r1 r2 1 1 PD – PB = 2T r1 r2 where h is the difference in water levels in the two limbs or PHYSICS Now, Given that 1 1 r1 r2 T = 0.07 Nm–1 , = 1000 kgm–3 3 3 3 mm = cm = m = 1.5 × 10–3 m, r2 = 3 × 10–3 m 20 20 100 2 2 0.07 1 1 m = 4.76 × 10–3 m = 4.76 mm 3 3 1000 9.8 1.5 10 3 10 Two narrow bores of diameters 3.0 mm and 6.0 mm are joined together to form a U-shaped tube open at both ends. If the U-tube contains water, what is the difference in its levels in the two limbs of the tube? Surface tension of water at the temperature of the experiment is 7.3 × 10–2 Nm–1. Take the angle of contact to be zero, and density of water to be 1.0 × 103 kg m–3 (g = 9.8 ms–2). Given that r1 = 3 .0 6.0 = 1.5 mm = 1.5 × 10–3 m, r2 = = 3.0 mm = 3.0 × 10–3 m, 2 2 T = 7.3 × 10–2 Nm–1, = 0º = 1.0 × 103 kg m–1, g = 9.8 ms–2 When angle of contact is zero degree, the radius of the meniscus equals radius of bore. 2T 2 7.3 10 2 Excess pressure in the first bore, P1 = r = = 97.3 Pascal 1.5 10 3 2 Excess pressure in the second bore, P2 = 2T 2 7.3 10 2 = = 48.7 Pascal r2 3 10 3 Hence, pressure difference in the two limbs of the tube P = P1 – P2 = hg or P1 P2 97.3 48.7 = = 5.0 mm g 1.0 10 3 9.8 CAPILLARY RISE IN A TUBE OF INSUFFICIENT LENGTH We know, the height through which a liquid rises in the capillary tube of radius r is given by 2T 2T h = Rg or h R = g = constant When the capillary tube is cut and its length is less then h (i.e. h'), then the liquid rises upto the top of the tube and spreads in such a way that the radius (R') of the liquid meniscus increases and it becomes more flat so that hR = h'R' = Constant. Hence the liquid does not overflow. r r If h' < h then R' > R or > cos ' cos cos ' < cos ' > Ex. 10 If a 5 cm long capillary tube with 0.1 mm internal diameter open at both ends is slightly dipped in water having surface tension 75 dyne cm–1, state whether (i) water will rise half way in the capillary. (ii) Water will rise up to the upper end of capillary (iii) Water will overflow out of the upper end of capillary/ Explain your answer. Sol. Given that surface tension of water, T = 75 dyne/cm 0.1 Radius r = mm = 0.05 mm = 0.005 cm, 2 density = 1 gm/cm3, angle of contact, = 0º. Let h be the height to which water rise in the capillary tube. Then 2T cos 2 75 cos 0º h= = cm = 30.58 cm. rg 0.005 1 981 But length of capillary tube, h’ = 5 cm h' (i) Because h > therefore the first possibility does not exist. 2 (ii) Because the tube is of insufficient length therefore the water will rise upto the upper end of the tube. (iii) The water will not overflow out of the upper end of the capillary. It will rise only upto the upper end of the capillary. PHYSICS The liquid meniscus will adjust its radius of curvature R’ in such a way that R’h’ = Rh 2T constant g where R is the radius of curvature that the liquid meniscus would possess if the capillary tube were of sufficient length Rh rh = h' h' r r r = 0.005 30.58 = 0.0306 cm cos cos 0º 5 APPLICATIONS OF SURFACE TENSION (i) The wetting property is made use of in detergents and waterproofing. When the detergent materials are added to liquids, the angle of contact decreases and hence the wettability increases. On the other hand, when water proofing material is added to a fabric, it increases the angle of contact, making the fabric water-repellant. The antiseptics have very low value of surface tension. The low value of surface tension prevents the formation of drops that may otherwise block the entrance to skin or a wound. Due to low surface tension the antiseptics spreads properly over the wound. The lubricating oils and paints also have low surface tension. So they can spread properly. Surface tension of all lubricating oils and paints is kept low so that they spread over a large area. Oil spreads over the surface of water because the surface tension of oil is less than the surface tension of cold water. A rough sea can be calmed by pouring oil on its surface. EFFECT OF TEMPERATURE AND IMPURITIES ON SURFACE TENSION The surface tension of a liquid decreases with the rise in temperature and vice versa. According to n Ferguson, T = T 0 1 where T 0 is surface tension at 0ºC, is absolute temperature of the liquid, c c is the critical temperature and n is a constant varies slightly from liquid and has mean value 1.21. This formula shows that the surface tension becomes zero at the critical temperature, where the interface between the liquid and its vapour disappears. It is for this reason that hot soup tastes better while machinery parts get jammed in winter. The surface tension of a liquid changes appreciably with addition of impurities. For example, surface tension of water increases with addition of highly soluble substances like NaCI, ZnSO 4 etc. On the other hand surface tension of water gets reduced with addition of sparingly soluble substances like phenol, soap etc. SURFACE ENERGY We know that the molecules on the liquid surface experience net downward force. So to bring a molecule from the interior of the liquid to the free surface, some work is required to be done against the intermolecular force of attraction, which will be stored as potential energy of the molecule on the surface. The potential energy of surface molecules per unit area of the surface is called surface energy. Unit of surface energy is erg cm –2 in C.G.S. system and Jm –2 in Sl system. Dimensional formula of surface energy is [ML°T –2 ] Surface energy depends on number of surfaces e.g. a liquid drop is having one liquid air surface while bubble is having two liquid air surface. RELATION BETWEEN SURFACE TENSION AND SURFACE ENERGY Consider a rectangular frame PQRS of wire, whose arm RS can slide on the arms PR and QS. If this frame is dipped in a soap solution, then a soap film is produced in the frame PQRS in fig. Due to surface tension (T), the film exerts a force on the frame (towards the interior of the film). Let be the length of the arm RS, then the force acting on the arm RS towards the film is. F = T × 2 [Since soap film has two surfaces, that is why the length is taken twice]. Let the arm RS be displaced to a new position R'S' through a distance x work done, W = Fx = 2Tx Increase in potential energy of the soap film. = EA = 2Ex = work done in increasing the area (W) where E = surface energy of the soap film per unit area. According the law of conservation of energy, the work done must be equal to the increase in the potential energy PHYSICS W A Thus, surface tension is numerically equal to surface energy or work done per unit increase surface area. 2T x = 2Ex or T = E = Ex. 11 A mercury drop of radius 1 cm is sprayed into 106 droplets of equal size. Calculate the energy expanded if surface tension of mercury is 35 × 10–3 N/m. Sol. If drop of radius R is sprayed into n droplets of equal radius r, then as a drop has only one surface, the initial surface area will be 4R2 while final area is n(4r2). So the increase in area S = n(4r2) – 4R2 So energy expended in the process, W = TS = 4T [nr2 – R2] .... (1) Now since the total volume of n droplets is the same as that of initial drop, i.e., 4/3(R3) = n[(4/3) r3] or r = R/n1/3 ....(2) Putting the value of r from equation (2) in (1) W = 4R2T (n)1/3 – 1]. Ex. 12 If a number of little droplets of water, each of radius r, coalesce to form a single drop of radius R, show that the rise in temperature will be given be 3T 1 1 J r R where T is the surface tension of water and J is the mechanical equivalent of heat. Sol. Let n be the number of little droplets. Since volume will remain constant, hence volume of n little droplets = volume of single drop 4 4 n × r3 = R3 or nr3 = R3 3 3 Decrease in surface area = n × 4r2 – 4R2 R3 1 1 R 2 = 4R3 r R r 1 1 Energy evolved W = T × decrease in surface area = T × 4R3 r R or nr 3 R 2 = 4 A = 4 [nr2 – R2] = 4 r W 4T R 3 1 1 = But Q = ms d r R J J where m is the mass of big drop, s is the specific heat of water and d is the rise in temperature. Heat produced, Q = 4T R 3 J 4 4T R 3 R 3 × 1 × 1 × d = 3 J 1 1 r R = volume of big drop × density of water × sp. heat of water × d 1 1 r R 3T 1 1 J r R Ex. 13 A film of water is formed between two straight parallel wires each 10 cm long and at a separation 0.5 cm. Calculate the work required to increase 1 mm distance between them. Surface tension of water = 72 × 10–3 N/m. 10cm Sol. Here the increase in area is shown by shaded portion in the figure. Since this is a water film, it has two surfaces, therefore 0.5cm increase in area, S = 2 × 10 × 0.1 = 2 cm2 Work required to be done, 0.1cm W = S × T = 2 × 10–4 × 72 × 10–3 = 144 × 10–7 joule = 1.44 × 10–5 joule
What Is Reflation? Reflation is a fiscal or monetary policy designed to expand output, stimulate spending, and curb the effects of deflation, which usually occurs after a period of economic uncertainty or a recession. The term may also be used to describe the first phase of economic recovery after a period of contraction. - Reflation is a policy that is enacted after a period of economic slowdown or contraction. - The goal is to expand output, stimulate spending and curb the effects of deflation. - Policies include tax cuts, infrastructure spending, increasing the money supply and lowering interest rates. Reflation aims to stop deflation—the general decline in prices for goods and services that occurs when inflation falls below 0%. It is a long-term shift, often characterized by a prolonged reacceleration in economic prosperity that strives to reduce any excess capacity in the labor market. Reflation policies typically include the following: - Reducing taxes: Paying lower taxes makes corporations and employees wealthier. It is hoped that extra earnings will be spent in the economy, lifting demand and prices for goods. - Lowering interest rates: Makes it cheaper to borrow money and less rewarding to stow capital away in savings accounts, encouraging people and businesses to spend more freely. - Changing the money supply: When central banks boost the amount of currency and other liquid instruments in the banking system the cost of money falls, generating more investment and putting more money in the hands of consumers. - Capital Projects: Large investment projects create jobs, boosting employment figures and the number of people with spending power. In short, reflationary measures aim to lift demand for goods by giving people and companies more money and motivation to spend more. Reflation policy has been used by American governments to try and restart failed business expansions since the early 1600s. Although almost every government tries in some form or another to avoid the collapse of an economy after a recent boom, none have ever succeeded in being able to avoid the contraction phase of the business cycle. Many academics believe government agitation only delays the recovery and worsens the effects. Example of Reflation In the wake of the Great Recession, the U.S. economy remained subdued and the Federal Reserve (FED) struggled to create inflation, even after utilizing a number of reflationary monetary policy tools, such as lower interest rates and increased money supply. It wasn't until the election of President Donald Trump that the economy got a sniff of fiscal reflation. President Trump pledged a trillion-dollar infrastructure bill and far-reaching tax cuts, hopeful that these measures would boost the economy to full capacity. Reflation vs. Inflation It is important not to confuse reflation with inflation. Firstly, reflation is not bad. It is a period of price increases when an economy is striving to achieve full employment and growth. Inflation, on the other hand, is often considered bad as it is characterized by rising prices during a period of full capacity. G.D.H. Cole once said, "reflation may be defined as inflation deliberately undertaken to relieve a depression." Additionally, prices rise gradually during a period of reflation and fast during a period of inflation. In essence, reflation can be described as controlled inflation.
PhD student discovers that solar radiation could be a more important source of lunar iron nanoparticles than previously thought. Tiny iron nanoparticles unlike any found naturally on Earth are nearly everywhere on the Moon—and scientists are trying to understand why. A new study led by Northern Arizona University doctoral candidate Christian J. Tai Udovicic, in collaboration with associate professor Christopher Edwards, both of NAU’s Department of Astronomy and Planetary Science, uncovered important clues to help understand the surprisingly active lunar surface. In an article recently published in Geophysical Research Letters, the scientists found that solar radiation could be a more important source of lunar iron nanoparticles than previously thought. Asteroid impacts and solar radiation affect the Moon in unique ways because it lacks the protective magnetic field and atmosphere that protect us here on Earth. Both asteroids and solar radiation break down lunar rocks and soil, forming iron nanoparticles (some smaller, some larger) that are detectable from instruments on satellites orbiting the Moon. The study used data from National Aeronautics and Space Administration (NASA) and Japan Aerospace Exploration Agency (JAXA) spacecraft to understand how quickly iron nanoparticles form on the Moon over time. “We have thought for a long time that the solar wind has a small effect on lunar surface evolution, when in fact it may be the most important process producing iron nanoparticles,” Tai Udovicic said. “Since iron absorbs a lot of light, very small amounts of these particles can be detected from very far away – making them a great indicator of change on the Moon.” Surprisingly, the smaller iron nanoparticles seemed to form at a similar rate as radiation damage in samples returned from the Apollo missions to the Moon, a hint that the Sun has a strong influence in their formation. “When I saw the Apollo sample data and our satellite data side by side for the first time, I was shocked,” Tai Udovicic said. “This study shows that the solar radiation could have a much larger influence in active change on the Moon than previously thought, not only darkening its surface, but it might also create small quantities of water usable in future missions.” As NASA prepares to land the first woman and the next man on the surface of the Moon by 2024 as part of the Artemis mission, understanding the solar radiation environment and possible resources on the Moon are critical. In future work recently awarded a NASA Future Investigators in Space Science and Technology (FINESST) grant, Tai Udovicic plans to broaden his targeted study to the entire Moon, but is also eager to take a closer look at mysterious lunar swirls, one of which was recently selected as a landing site for the upcoming Lunar Vertex rover. He also studies lunar temperatures and water ice stability to inform future missions. “This work helps us understand, from a bird’s eye view, how the lunar surface changes over time,” said Tai Udovicic. “While there is still a lot to learn, we want to make sure that when we have boots back on the Moon, that those missions are backed by the best science available. It’s the most exciting time to be a lunar scientist since the tail end of the Apollo era in the 70s.” Reference: “New Constraints on the Lunar Optical Space Weathering Rate” by C. J. Tai Udovicic, E. S. Costello, R. R. Ghent and C. S. Edwards,19 June 2021 , Geophysical Research Letters.
How do you change those funny numbers and letters to something you or your computer can understand? Converting hexadecimal to binary is very easy, which is why hexadecimal has been adopted in some programming languages. Converting to decimal is a little more involved, but once you've got it it's easy to repeat for any number. Part 1 of 3: Converting Hexadecimal to Binary 1Convert each hexadecimal digit to four binary digits. Hexadecimal was adopted in the first place because it's so easy to convert between the two. Essentially, hexadecimal is used as a way to display binary information in a shorter string. This chart is all you need to convert from one to the other: X Research source Hexadecimal Binary 0 2Try it yourself. It really is as simple as changing the digit into the four equivalent binary digits. Here are a few hex numbers for you to convert. Highlight the invisible text to the right of the equal sign to check your work: - A23 = 1010 0010 0011 - BEE = 1011 1110 1110 - 70C558 = 0111 0000 1100 0101 0101 1000 3Understand why this works. In the "base two" binary system, n binary digits can be used to represent 2n different numbers. For example, with four binary digits, you can represent 24 = 16 different numbers. Since hexadecimal is a base sixteen system, a one digit number can be used to represent 161 = 16 different numbers. This makes conversion between the two systems extremely easy. X Research source - You can also think of this as the counting systems "flipping over" to another digit at the same time. Hexadecimal counts "...D, E, F, 10" at the same time binary counts "1101, 1110, 1111, 10000". Part 2 of 3: Converting Hexadecimal to Decimal 1Review how base ten works. You use decimal notation every day without having to stop and think about the meaning, but when you first learned it, your parent or teacher might have explained it to you in more detail. A quick review of how ordinary numbers are written will help you convert the number: X Research source - Each digit in a decimal number is in a certain "place." Moving from right to left, there's the "ones place," "tens place," "hundreds place," and so on. The digit 3 just means 3 if it's in the ones place, but it represents 30 when located in the tens place, and 300 in the hundreds place. - To put it mathematically, the "places" represent 100, 101, 102, and so on. This is why this system is called "base ten," or "decimal" after the Latin word for "tenth." 2Write a decimal number as an addition problem. This will probably seem obvious, but it's the same process we'll use to convert a hexadecimal number, so it's a good starting point. Let's rewrite the number 480,13710. (Remember, the subscript 10 tells us the number is written in base ten.): - Starting with the rightmost digit, 7 = 7 x 100, or 7 x 1 - Moving left, 3 = 3 x 101, or 3 x 10 - Repeating for all digits, we get 480,137 = 4x100,000 + 8x10,000 + 0x1,000 + 1x100 + 3x10 + 7x1. 3Write the place values next to a hexadecimal number. Since hexadecimal is base sixteen, the "place values" correspond to the powers of sixteen. To convert to decimal, multiply each place value by the corresponding power of sixteen. Start this process by writing the powers of sixteen next to the digits of a hexadecimal number. We'll do this for the hexadecimal number C92116. Start on the right with 160, and increase the exponent each time you move left to the next digit: X Research source - 116 = 1 x 160 = 1 x 1 (All numbers are in decimal except where noted.) - 216 = 2 x 161 = 2 x 16 - 916 = 9 x 162 = 9 x 256 - C = C x 163 = C x 4096 4Convert alphabetic characters to decimal. Numerical digits are the same in decimal or hexadecimal, so you don't need to change them (for instance, 716 = 710). For alphabetic characters, refer to this list to change them to the decimal equivalent: - A = 10 - B = 11 - C = 12 (We'll use this on our example from above.) - D = 13 - E = 14 - F = 15 5Perform the calculation. Now that everything is written in decimal, perform each multiplication problem and add the results together. A calculator will be handy for most hexadecimal numbers. Continuing our example from earlier, here's C921 rewritten as a decimal formula and solved: X Research source - C92116 = (in decimal) (1 x 1) + (2 x 16) + (9 x 256) + (12 x 4096) - = 1 + 32 + 2,304 + 49,152. - = 51,48910. The decimal version will usually have more digits than the hexadecimal version, since hexadecimal can store more information per digit. 6Practice the conversion. Here are a few numbers to convert from hexadecimal into decimal. Once you've worked out the answer, highlight the invisible text to the right of the equal sign to check your work: - 3AB16 = 93910 - A1A116 = 4137710 - 500016 = 2048010 - 500D16 = 2049310 - 18A2F16 = 10091110 Part 3 of 3: Understanding Hexadecimal Basics 1Know how to use hexadecimal. Our ordinary decimal counting system is base ten, using ten different symbols to display numbers. Hexadecimal is a base sixteen number system, meaning it uses sixteen characters to display numbers. X Research source - Counting from zero upward: Hexadecimal Decimal Hexadecimal Decimal 0 - Counting from zero upward: 2Use subscript to show which system you're using. Whenever it might be unclear which system you're using, use a decimal subscript number to denote the base. For example, 1710 means "17 in base ten" (an ordinary decimal number). 1710 = 1116, or "11 in base sixteen" (hexadecimal). You can skip this if your number has an alphabetic character in it, such as B or E. No one will mistake that for a decimal number.Advertisement QuestionWhat is 480137 converted to a hexadecimal and then to a binary number?Community Answer480137 converted to a hexadecimal is 75389, and the binary is 0111 0101 0011 1000 1001. QuestionWhat is the decimal equivalent of the hex number 0x3F?Community Answer63. You just need to convert the 3F part. QuestionHow many permutations are there when using the hexadecimal system?Community AnswerIt depends on how many hexadecimal digits you are using. 1 digit hexadecimal has 16 permutations. 2 digits have 256 permutations, etc. Each hexadecimal digit is represented with 4 bits, so if you have x hexadecimal digits, there are 2 ^ (x * 4) possible permutations. QuestionHow do I convert binary numbers to hexadecimal?Community AnswerGroup the bits in terms of 4, and calculate the decimal value of each group. If you get values above 9, like 10, it's A, 11 is B.....15 is F, etc. QuestionHow do I change 38 hexadecimal to a binary number?Community AnswerFirst, let's turn it into a denary number, just to make things easier: 3 x 16 = 48, 1 x 8 = 8; 48 + 8 = 56. Now turn our denary value into a binary one: 1 goes first in the 32 column, 56 - 32 leaves 24. One goes in the 16 column, 24 - 16 leaves 8. The remaining 1 goes in the 8 column; the binary result is 0011 1000. QuestionHow do I change 18A into binary?Top AnswererUse the method specified in the article: Hexadecimal 1 is binary 0001, hexadecimal 8 is binary 1000, hexadecimal A is binary 1010, so the number 18A in binary is 110001010. This works because 2^4 = 16. This is also more detailed explained in the article itself. QuestionHow do I convert hexadecimal numbers to binary?Community AnswerConvert each digit to binary, then add them. For example, 18A =110001010 as if you divide them into 4 groups. You will find that the first group at the right is equivalent to A and the second one is equivalent to 8 and the same with 1. What is 102632 converted to a hexadecimal and a binary number? - Long hexadecimal numbers may require an online calculator to convert to decimal. You can also skip the work and have an online converter do the work for you, although it's a good idea to understand how the process works. X Research source - You can adapt the "hexadecimal to decimal" conversion to convert any other base x numbering system to decimal. Just replace the powers of sixteen with powers of x instead. Try learning the base-60 Babylonian counting system! X Research source - ↑ https://www.swarthmore.edu/NatSci/echeeve1/Ref/BinaryMath/NumSys.html - ↑ https://kb.iu.edu/d/agxz - ↑ http://www.uwosh.edu/coehs/cmagproject/concepts/documents/Developing_Base_Ten_Understanding.pdf - ↑ http://csc.columbusstate.edu/woolbright/CONV.HTM - ↑ http://www.sci.brooklyn.cuny.edu/~jones/cisc1110/basesystems.pdf - ↑ http://homepage.smc.edu/morgan_david/cs40/hex-system.htm - ↑ http://www.mathsisfun.com/binary-decimal-hexadecimal-converter.html - ↑ http://www-history.mcs.st-and.ac.uk/HistTopics/Babylonian_numerals.html About This Article To convert hexadecimal to binary, convert each hexadecimal digit to 4 binary digits. Each of the 16 hexadecimal digits are equal to 4 binary digits, so all you need to do is memorize the 16 conversions. For example, the hexadecimal 1 is equal to the binary 0001. To convert hexadecimal to decimal, multiply each place value in the hexadecimal number by the corresponding power of sixteen. Then, add all of the products together to get the decimal. If you want to learn how to label the systems you're using in subscripts, keep reading the article!
Bacterial DNA Replication and Cell Division Help (page 2) Bacterial DNA Replication The circular chromosome of bacteria presents special problems for replication. Circular chromosomes usually have a single site, called the origin, or ori, site, at which replication originates. By contrast, many ori sites exist on each chromosome of eukaryotes. Once the replication process starts, it usually proceeds bidirectionally from the ori site to form two replication forks. As the two strands of a right-handed, double-helical, circular DNA unwind during replication, the molecule tends to become positively supercoiled or overwound, i.e., twisted in the same direction as the strands of the double helix. These supercoils are so tight that they would interfere with further replication if they were not removed. Topoisomerases are a group of enzymes that can change the topological or configurational shape of DNA. DNA gyrase is a bacterial topoisomerase that makes double-stranded cuts in the DNA, holds on to the broken ends so they cannot rotate, passes an intact segment of DNA through the break, and then reseals the break on the other side (Fig. 10-2). Fig. 10-2. A proposed mechanism whereby DNA gyrase "pumps" negative supercoiling into DNA. A relaxed, covalently closed, circular DNA molecule (a) is bent into a configuration for strand passage (b). DNA gyrase makes double-stranded cuts (c), holds on to the ends, passes an intact segment through the break, and reseals the break on the other side (d). This action of DNA gyrase quickly removes positive supercoils and momentarily relaxes the DNA molecule into a more energetically stable state. However, with the expenditure of energy, DNA gyrase normally pumps negative supercoiling or under winding (twisting in a direction opposite to the turns of the double helix) into relaxed DNA circles so that virtually all DNAs in both prokaryotes and eukaryotes naturally exist in the negative super-coiled state. Relaxed circles and positively super-coiled DNA exist only in the laboratory. Localized regions of DNA transiently and spontaneously unwind to single-stranded "bubbles" and then return to their former topology as hydrogen bonds between complementary base pairs are broken and reformed by thermal agitation. The strain of underwinding is thus momentarily relieved in a superhelix by an increase in the number, size, and duration of these bubbles. An equilibrium normally exists between these super-coiled and "bubbled" states. More bubbles form as the temperature increases. At each replication fork, an enzyme called helicase unwinds the two DNA strands. Single-stranded, DNA binding (SSB) proteins protect the single-stranded regions in the replication forks from forming intrastrand base pairings that could cause a tangle of partially double-stranded segments that would interfere with replication. The enzyme primase synthesizes short RNA primers using a region on each strand as a template. Primers are required for DNA polymerase to begin extending the new DNA strand because DNA polymerase requires a 3'OH to initiate the bonding reaction. Three DNA polymerase enzymes (referred to as pol I, pol II, and pol III) have been found in E. coli. Pol III is the principal replicating enzyme. Gaps left by pol III are filled by pol I, and DNA ligase seals the nicks. The function of pol II is not well established, although it is known that it is not involved in RNA primer replacement. In addition to their 5' to 3' synthetic activity, both pol I and pol III have 3' to 5' exonuclease activity, which plays a "proofreading" role by removing mismatched bases mistakenly inserted during chain polymerization. Pol I also has 5' to 3' exonuclease activity by which it normally removes primers and replaces them with complementary DNA sequences after polymerization has begun. About halfway through the above replication process, the replicative intermediate molecule looks like the Greek letter theta (θ), so is referred to as theta replication (Fig. 10-3). Another type of bacterial replication is used to transfer a linear DNA molecule during bacterial conjugation or for the production of linear phage genomes. A nick occurs in one strand of a DNA double helix, creating free 3'-OH and 5'-P termini. Helicase and SSB proteins establish a replication fork. No primer is necessary because a strand with a free 3'-OH is available for elongation by DNA polymerase III as the leading strand. Simultaneously with replication of the leading strand, the template for the lagging strand is displaced. The displaced strand is discontinuously replicated to produce Okazaki fragments in the usual way (see Fig. 3-11). The result of this replication model is a circle with a linear tail, resembling the Greek letter sigma (σ). Hence, this model is called sigma replication or rolling-circle replication (Fig. 10-4). The circle may revolve several times, creating concatemers or covalently connected, linear repetitions of bacterial genomes. An endonuclease makes cuts at slightly different positions on each DNA strand of the concatemer to create genome-sized segments containing "sticky ends" (single-stranded complementary ends). The linear genomes circularize by base pairing of the sticky ends. DNA ligase seals each gap to create covalently closed (circular), double-stranded DNA molecules. A replicating bacterial chromosome is thought to be attached to invaginations of the cell membrane at each replication fork. After DNA replication, the cell elongates by growth of the sector between the two attachment points, causing the two chromosomal replicas to move apart. A septum of new cell membrane and wall is then synthesized between the two chromosomes, creating two progeny cells (Fig. 10-5). The "passing down" of DNA from parent to progeny cell is called vertical gene transfer and the overall process of bacterial cell division is called binary fission. Fig. 10-5. A model for segregation of bacterial DNA ("chromosome") replicas. (1) A circular DNA molecule is attached to invaginations of the cell membrane at two points. The DNA is a theta structure (about half replicated). (2) Replication is complete. (3) Cell division is beginning via growth of the membrane region shown in medium gray. Both daughter chromosomes are already partially replicated. (4) Cell division is complete. (5) Daughter cells are midway through the next generation. The membrane attachment points have moved to the center of the cell as a result of growth of new membrane (dark gray) after cell division. (After G. S. Stent and R. Calendar, Molecular Genetics, 2nd ed., W. H. Freeman and Company, New York, 1978.) When bacteria are growing exponentially, most cells contain two to four identical chromosomes in various stages of replication. If a mutation occurs during replication, the new copy of DNA will be slightly different from the parental template. If this mutation occurs in a coding region that results in a defective protein, the effect of the mutation (i.e., the phenotype) will only be observable once the new cell has depleted its level of wild-type protein. This is a phenomenon known as phenotypic lag. EXAMPLE 10.4 Resistance to a specific bacteriophage can be acquired by mutation of a gene responsible for the phage receptor on the cell's surface. Resistance cannot be fully realized until the receptor sites (synthesized under the direction of the former phage-sensitive genotype) have been completely diluted out through successive cell divisions. If even one receptor remains on the mutant cell, it is still susceptible to phage infection. Thus, many cell generations may be required before a phage-resistant mutation can be fully expressed in a progeny cell. Practice problems for these concepts can be found at: - Kindergarten Sight Words List - First Grade Sight Words List - 10 Fun Activities for Children with Autism - Signs Your Child Might Have Asperger's Syndrome - A Teacher's Guide to Differentiating Instruction - Theories of Learning - Child Development Theories - Social Cognitive Theory - Curriculum Definition - Why is Play Important? Social and Emotional Development, Physical Development, Creative Development
The resulting discovery of dark energy and the accelerating universe rewrote our understanding of the cosmos. Yet the origin of these supernovae, which have proved so useful, remains unknown. "The question of what causes a Type Ia supernova is one of the great unsolved mysteries in astronomy," says Rosanne Di Stefano of the Harvard-Smithsonian Center for Astrophysics (CfA). Astronomers have very strong evidence that Type Ia supernovae come from exploding stellar remnants called white dwarfs. To detonate, the white dwarf must gain mass until it reaches a tipping point and can no longer support itself. There are two leading scenarios for the intermediate step from stable white dwarf to supernova, both of which require a companion star. In the first possibility, a white dwarf swallows gas blowing from a neighboring giant star. In the second possibility, two white dwarfs collide and merge. To establish which option is correct (or at least more common), astronomers look for evidence of these binary systems. Given the average rate of supernovae, scientists can estimate how many pre-supernova white dwarfs should exist in a galaxy. But the search for these progenitors has turned up mostly empty-handed. To hunt for accreting white dwarfs, astronomers looked for X-rays of a particular energy, produced when gas hitting the star's surface undergoes nuclear fusion. A typical galaxy should contain hundreds of such "super-soft" X-ray sources. Instead we see only a handful. As a result, a recent paper suggested that the alternative, merger scenario was the source of Type Ia supernovae, at least in many galaxies. That conclusion relies on the assumption that accreting white dwarfs will appear as super-soft X-ray sources when the incoming matter experiences nuclear fusion. Di Stefano and her colleagues have argued that the data do not support this hypothesis. In a new paper, Di Stefano takes the work a step further. She points out that a merger-induced supernova would also be preceded by an epoch during which a white dwarf accretes matter that should undergo nuclear fusion. White dwarfs are produced when stars age, and different stars age at different rates. Any close double white-dwarf system will pass through a phase in which the first-formed white dwarf gains and burns matter from its slower-aging companion. If these white dwarfs produce X-rays, then we should find roughly a hundred times as many super-soft X-ray sources as we do. Since both scenarios - an accretion-driven explosion and a merger-driven explosion - involve accretion and fusion at some point, the lack of super-soft X-ray sources would seem to rule out both types of progenitor. The alternative proposed by Di Stefano is that the white dwarfs are not luminous at X-ray wavelengths for long stretches of time. Perhaps material surrounding a white dwarf can absorb X-rays, or accreting white dwarfs might emit most of their energy at other wavelengths. If this is the correct explanation, says Di Stefano, "we must devise new methods to search for the elusive progenitors of Type Ia supernovae." Di Stefano's paper has been accepted for publication in The Astrophysical Journal and is available online.Headquartered in Cambridge, Mass., the Harvard-Smithsonian Center for Astrophysics (CfA) is a joint collaboration between the Smithsonian Astrophysical Observatory and the Harvard College Observatory. CfA scientists, organized into six research divisions, study the origin, evolution and ultimate fate of the universe. For more information, contact:David A. Aguilar Christine Pulliam \| EurekAlert! SF State astronomer searches for signs of life on Wolf 1061 exoplanet 20.01.2017 \| San Francisco State University Molecule flash mob 19.01.2017 \| Technische Universität Wien An important step towards a completely new experimental access to quantum physics has been made at University of Konstanz. The team of scientists headed by... Yersiniae cause severe intestinal infections. Studies using Yersinia pseudotuberculosis as a model organism aim to elucidate the infection mechanisms of these... Researchers from the University of Hamburg in Germany, in collaboration with colleagues from the University of Aarhus in Denmark, have synthesized a new superconducting material by growing a few layers of an antiferromagnetic transition-metal chalcogenide on a bismuth-based topological insulator, both being non-superconducting materials. While superconductivity and magnetism are generally believed to be mutually exclusive, surprisingly, in this new material, superconducting correlations... Laser-driving of semimetals allows creating novel quasiparticle states within condensed matter systems and switching between different states on ultrafast time scales Studying properties of fundamental particles in condensed matter systems is a promising approach to quantum field theory. Quasiparticles offer the opportunity... Among the general public, solar thermal energy is currently associated with dark blue, rectangular collectors on building roofs. Technologies are needed for aesthetically high quality architecture which offer the architect more room for manoeuvre when it comes to low- and plus-energy buildings. With the “ArKol” project, researchers at Fraunhofer ISE together with partners are currently developing two façade collectors for solar thermal energy generation, which permit a high degree of design flexibility: a strip collector for opaque façade sections and a solar thermal blind for transparent sections. The current state of the two developments will be presented at the BAU 2017 trade fair. As part of the “ArKol – development of architecturally highly integrated façade collectors with heat pipes” project, Fraunhofer ISE together with its partners... 19.01.2017 \| Event News 10.01.2017 \| Event News 09.01.2017 \| Event News 20.01.2017 \| Awards Funding 20.01.2017 \| Materials Sciences 20.01.2017 \| Life Sciences
The Focault Pendulum is a famous experiment which is alleged to give simple, direct evidence of the earth's rotation. Introduced in the 1851, Léon Foucault claimed that the swinging rotational motions were proof the earth's rotation. The pendulum swings back and forth, and over time, slowly seems to rotate over its arena or "pit". It is explained that the earth is rotating beneath the pendulum. Today Foucault Pendulums are popular displays in science museums and universities. A common criticism of the Focault Pendulum is to point out that when the pendulum experiment has been recreated and put into motion, the pendulum has often been seen to rotate in excess, in shortness, or in an opposite direction from the direction it should have traveled according to theory. At times it does not rotate at all. Those scientists who have repeated the experiment have freely admitted that “it was difficult to avoid giving the pendulum some slight lateral bias at starting.” In the unmotorized Focault Pendulum experiment the pendulum is, as we will read below, generally inconsistent in its movements. Because of air resistance and the impossibility of perpetual motion, the unmotorized pendulum will only move for a short while before needing to be reset. In motorized Focault Pendulums, as seen in museum exhibits, it is the repetitive machinery which imparts the repeating lateral bias that creates the regular results seen for the museum's visitors. Thus, the experiment is entirely invalid as a demonstration of diurnal rotation. That a pendulum on a line can rotate as it swings back and fourth has more to do with the initial conditions which set it into motion than the supposed rotation of the earth. Lady Blount provides the following in The Romance of Science: “ This pendulum, modern scientists tell us, affords a visible proof that we are living on a whirling globe, which, according to a ‘Work on Science’ now before me, is spinning upon its so-called axis at the rate of over 1,000 miles an hour at the equator; and, in addition to other motions, is rushing on an everlasting tour round the sun (the diameter of which is said to be 813,000 miles, and its weight 354,936 times greater than the earth from which it is said to be about 93,000,000 miles distant,) at the rate of over 1,000 miles per minute. Now to prove that the earth really has these motions a pendulum is suspended at the show; the showman sets motion, and bids the gaping world of thoughtless men and women to ‘behold a proof’ that we are living on a whirling globe which is rushing away through space! ” In Earth Not a Globe Samuel Birley Rowbotham informs us that the variation of the pendulum is often non-uniform and unpredictable: “ First, when a pendulum, constructed according to the plan of M. Foucault, is allowed to vibrate, its plane of vibration is often variable – not always. The variation when it does occur, is not uniform – is not always the same in the same place; nor always the same either in its rate or velocity, or in its direction. It cannot therefore be taken as evidence; for that which is inconstant cannot be used in favor of or against any given proposition. It therefore is not evidence and proves nothing! Secondly, if the plane of vibration is observed to change, where is the connection between such change and the supposed motion of the Earth? What principle of reasoning guides the experimenter to the conclusion that it is the Earth which moves underneath the pendulum, and not the pendulum which moves over the Earth? What logical right or necessity forces one conclusion in preference to the other? Thirdly, why was not the peculiar arrangement of the point of suspension of the pendulum specially considered, in regard to its possible influence upon the plane of oscillation? Was it not known, or was it overlooked, or was it, in the climax of theoretical revelry, ignored that a ‘ball-and-socket’ joint is one which facilitates circular motion more readily than any other? ” The Wrong Direction The Focault Pendulum is often seen to move in the wrong direction entirely. See the following from A Hundred Proofs The Earth is Not a Globe: “ Astronomers have made experiments with pendulums which have been suspended from the interior of high buildings, and have exulted over the idea of being able to prove the rotation of the Earth on its ‘axis,’ by the varying direction taken by the pendulum over a prepared table underneath – asserting that the table moved round under the pendulum, instead of the pendulum shifting and oscillating in different directions over the table! But, since it has been found that, as often as not, the pendulum went round the wrong way for the ‘rotation’ theory, chagrin has taken the place of exultation, and we have a proof of the failure of astronomers in their efforts to substantiate their theory. ” From The Romance of Science (8-10) we read: “ We believe, with all due deference to the pendulum, and its proprietor, that it proves nothing but the craftiness of the inventor; and we can only describe the show and showman as deceptions. A thing so childish as this ‘pendulum proof’ that it can only be described as one of the most simple and ridiculous attempts to gull the public that has ever been conceived. …It has been said that the pendulum experiment proves the rotation of the earth, but this is quite impossible, for one pendulum turns one way; and sometimes, another pendulum turns in the opposite direction. Now we ask does the earth rotate in opposite directions at different places at one and the same time? We should like to know. Perhaps the experimenters will kindly enlighten us on this point. …If the earth had the terrible motions attributed to it, there would be some sensible effects of such motions. But we neither feel the motion, see it, nor hear it. And how people can stand watching the pendulum vibrate, and think that they are seeing a proof of the motions of the earth, almost passes comprehension. They are, however, brought up to believe it, and it is thought to be ‘scientific’ to believe what the astronomers teach. ” Museum Exhibit Devices We read the following from an installation guide: “ Pay close attention to the photo beams alignment. This adjustment can effect the Ball’s precession around the pit. It may require a couple of days to determine if precession is operating properly. Precession is a function of the Earths rotation. ” We are instructed to spend several days adjusting the alignment of the photo beams, which affects the pendulum's precession, an element which is supposedly a function of the earth's rotation, until we have determined that the "precession is operating properly". South Pole Pendulum In 2001 a Foucault pendulum device was installed at the "South Pole" in the stairwell of a newly constructed research station. It has been claimed that this verified Foucault’s theory. However, excerpts from the report show that adjustments had to be made to obtain the desired results. “ A hole was drilled through a length of 2×4 through which the pendulum’s wire was fixed in such a way that it would have freedom of movement in all directions. The 2×4 was then fixed to the top of the stairwell and the wire suspended down into the stairwell. ” “ There was no mechanism for keeping the pendulum swinging and the amplitude decayed within a couple of hours so it had to be restarted periodically over the 24 hour period. ...Our first attempt with the pendulum showed the Earth spinning backward from what was expected. We didn’t notice this at first because we’re all from the Northern Hemisphere and are accustomed to the earth spinning in an anticlockwise direction. We then realized that from our frame of reference the earth should be spinning clockwise so we had to modify the pendulum. ...Our second attempt showed the earth rotating in the proper direction but at an angular velocity twice what is expected (ire., 12 hour days instead of 24). We suspected some kind of government conspiracy but decided to make a further modification and try it again. ...It was difficult to make the pendulum swing in a plane instead of an ellipse. After several attempts with various techniques of holding the bob and dropping it we always got some kind of ellipse instead of a plane. This adds to our error because it is more difficult to locate and mark the pendulum arc’s apex. A way to do it is to suspended the bob by tying it off with a piece of string and letting it settle, then burn through the string. The bob would then drop without any outside force and swing in a plane. Since it is against the Antarctica Treaty to have any open flames at the South Pole we could not do this. After much practice Mike Town got very adept at dropping the bob so that it arced in a plane. ” One alternative explanation that has been suggested by those who do accept the Focault Pendulum is Mach's Principle. Mach's Principle explains that if the earth was still and the all the stars went around the Earth then the gravitational pull of the stars would pull the pendulum. As Mach said "The universe is not twice given, with an earth at rest and an earth in motion; but only once, with its relative motions alone determinable. It is accordingly, not permitted us to say how things would be if the earth did not rotate."
Polynomial Teacher Resources Find Polynomial educational ideas and activities Showing 1 - 20 of 1,379 resources This is an interesting problem that approaches higher degree polynomials from a different perspective. Given a third degree polynomial, two known zeros, and a y-intercept, find the value of the polynomialÕs coefficients. This problem challenges the learner to gain a deeper understanding of zeros, graphs of polynomials, factors, and the algebra behind multiplying polynomials. This problem has a number of correct solutions, but that is what makes it challenging and still within the reach of a second-year algebra student. Though it "sounds like a really fancy word," polynomials prove to be no match for Sal's mathematical skills. After defining the term and providing a few examples, Sal works through a few equations that add or subtract polynomials, showing how variable or constant terms that are raised to non-zero exponents can be simple. Put your math pupils’ division skills - and algebra skills - to the test with this video, which introduces the concept of polynomial division. Sal's examples increase in complexity as the video progresses, allowing viewers to see how basic skills can be applied to both simple and complicated problems. This worksheet first reviews the definition of a polynomial and a like term. Then learners use the distributive property to rewrite expressions where polynomials are multiplied by monomials. Finally they are introduced to both the FOIL method and box model for multiplying binomials and trinomials. This lesson is all about polynomial functions, graphs, and parabolas! It even comes with quite a few follow along worksheets. Get your class graphing polynomial functions with a degree higher than two, identifying the function represented by a graph, and finding the zeros and intercepts. A very thorough lesson. Explore the concept of graphing polynomials with your class. Scholars graph polynomials and determine their end behavior. They use their calculator to determine the end behavior of linear, quadratic, and cubic equations. This comprehensive lesson looks at higher-degree polynomials. The important topics covered include the Remainder Theorem, the Factor Theorem, and synthetic division. Included is a clearly written tutorial, problem sets, summary, and answer key. Define, simplify, add, and subtract polynomials. Identify the degrees of a polynomial. Determine if a polynomial is a monomial, a binomial, or a trinomial. Then look at simplifying a polynomial by combining the coefficient values of like terms. This resource contains clear steps leading up to and including how to add and subtract polynomials. Note: Be careful with subtracting polynomials because the signs of the terms inside the parenthesis will change. This algebra worksheet reviews writing polynomials in standard form from factored form, looks at the graphs of polynomials of degree higher than two, and identifies the zeros of polynomials using the Factor Theorem and Fundamental Theorem of Algebra. The instructor defines a prime polynomial as much the same as a prime number: they only have two factors, 1 and itself. She further demonstrates what a composite polynomial is to give an example of a polynomial that is not prime. A polynomial is the sum of one or more monomials. They are identified by the number of terms they have. Monomials, binomials, trinomials, find out what these terms identify. There are four terms in this polynomial. Try to factor it by grouping the terms. Confused? Try to get four terms into the product of two binomials? Still confused? Watch the instructor as she goes through the steps to group and regroup terms to get the expression into two binomials. Don't forget to check the solution. So you want to learn how to add polynomials? Here are the steps: rewrite by removing the parentheses, identify the like terms, group like terms together, combine the like terms, and finally rewrite the simplified expression. In this Algebra II/college level worksheet, students use the discriminate to determine if a polynomial has factors with integral coefficients and factor polynomials, including factoring by pulling out the GCF and factoring by grouping. The two page worksheet contains seventy-two problems. Answers are included. It's not a complicated word problem, but the picture frame just gives the problem a little context. So you are trying to find the difference, which means subtraction. The instructor reviews how to distribute the negative sign to the second polynomial. Then it's just a matter of combining like terms. This clip is great! Your budding algebra experts will have no trouble understanding what the degree in a polynomial means. The degree is compared to a mountain range, the highest peak in the range usually names the entire set of mountains. This is a well constructed and very helpful way to help struggling learners understand polynomials. Learn how to factor the quadratic polynomial ax(squared)+bx+c, where (a) does not equal 1. This quick algebra lesson focuses on using the "bottoms up" method to factor this polynomial. The instructor uses an interactive whiteboard to show an example. Explore finding zeros of polynomial functions in a number of ways including using synthetic division, depressed polynomials, factoring, the Factor Theorem. and the Rational Roots Theorem. This lesson also looks at complex conjugates as zeros of a polynomial function. Help your pupils define a Taylor polynomial approximation to a function f of degree n about a point x = a. After completing several problems with guided practice, individuals graph convergence of Taylor polynomials and use them to approximate function values. High schoolers generate the equation of a polynomial given its roots and the end behavior of the function. They need to apply theorems concerning the multiplicity of roots, conjugates of irrational or imaginary roots to find a polynomial. Additionally, they will use a graphing utility to determine local extrema.
Are you curious about DNS (domain name system), the system responsible for routing traffic for every domain on the web? In this post, you will learn what DNS is and exactly how it works. Let’s get started: What is DNS? The domain name system (DNS) allows us to access websites with an alphanumeric web address. The world wide web as we know it was invented in 1989, and the first web page didn’t go online until 1991. Still, the internet was being developed and in use decades before then. Your website and other entities hosted on the web have a specific location on the net. This is represented by a numeric IP address, such as 220.127.116.11, similar to how your street address represents the location of your home. Domains, such as bloggingwizard.com, didn’t exist when the internet was being developed. Its users were required to enter a location’s IP address instead. Given how difficult it is to remember and enter numeric addresses for everything the internet’s users needed to access at this time, finding a new method to access them was crucial. Paul Mockapetris brought this new method to the internet when he invented the domain name system in 1983. By 1984, internet users could access the net’s locations with user-friendly, alphanumeric domain names and six top level domains (TLDs): - .com – created for commercial purposes. - .org – created for organizations. - .net – created for networks. - .gov – created for government-sponsored locations. - .edu – created for educational computer systems. - .mil – created for military-sponsored locations. DNS is responsible for translating every domain on the web into an identifiable IP address. When you want to visit a website, your browser relies on this system to find its exact location on the web. What is a nameserver? You may find that some people use the terms DNS and nameservers interchangeably. Typically because they are referencing the exact same thing – your DNS records. In a technical sense, a nameserver is simply the server where DNS records are stored. How do domains work? In order to understand how domain servers work, we need to understand how domains work. Domains are the alphanumeric addresses we use to access specific locations on the web, typically websites. As we explained earlier, they represent the IP addresses that identify those locations and allow us to access them without having to enter those IP addresses into our address bars. When you enter a web address into your browser, the DNS goes through several steps before your browser loads the web page you’re trying to access a second later. In order for your browser to complete your request, it must receive the IP address of the domain you’re trying to access from the DNS. This is called DNS resolution, and it runs through a few different nameservers before it completes your request. This includes the TLD nameserver. TLD stands for “top level domain.” Domains have a hierarchy made up of three levels, though modern domains only use the second and top levels. Here’s an example featuring the domain of the tool I’m using to write this draft, Google Docs. Google Docs’ domain – docs.google.com: - docs = third level or “subdomain.” - .google = second level or “domain name.” - .com = top level or “domain extension.” Remember when we said there were only six TLDs in 1984? Today, there are more than 1,500. They’re organized into three different categories. Generic top level domains (gTLD) are the biggest category. gTLDs include common domains like .com, .org and .net but also include more unique iterations. When you register a domain today, you’ll find offers for domains that include such TLDs as .biz, .me, .io, .xyz, .pizza, .beer, .motorcycles and more. Sponsored top level domains (sTLD) are TLDs sponsored by specific entities, such as governments, military forces and educational organizations. As such, these TLDs include .gov, .mil and .edu. Country code top level domains (ccTLD) are TLDs made for specific countries. Websites use them when they want to target customers in specific countries. There are more than 200 ccTLDs in existence, including .uk for the United Kingdom, .ru for Russia, .cn for China, .br for Brazil, and so on and so forth. When you register a domain, you need to choose a domain name and TLD for it. Its IP address will be stored on your registrar’s DNS server. It’s important to note that you will not have authority over other domains that use your domain name with different TLDs unless you register it. This means if you register example.com, a competitor could register example.xyz. They’re treated as entirely different domains by the DNS. In order to have your new domain lead to your website when you enter it in your browser, you must use your registrar’s DNS settings to point the domain to your host’s nameservers. How do domain servers work? Nameservers are part of the process involved in translating domains into their loctable IP addresses. They store DNS records, particularly those very IP addresses that help us identify websites. Let’s go over the process (called DNS resolution) the DNS goes through to return an IP address to your browser when you try to visit a website. Let’s say you want to visit the Google Docs dashboard. You enter “docs.google.com” in your browser (or your browser does if you use a shortcut). Before the DNS can translate that domain for you, it needs to run your request through four primary servers in order to identify its IP address. The first is the recursor server. This one is simple as its purpose is to simply handle your request. It’ll also send additional requests for you if need be. Next is the root nameserver. Nameservers are containers for DNS records, including the A record that contains a domain’s IP address. We’ve established this already. We’ve also already established how the DNS is responsible for translating human-readable domains into machine-friendly IP addresses through a process called DNS resolution. The root nameserver initiates this process. After your request moves through the root nameserver, it moves onto the TLD nameserver. At this point, the DNS is looking for your domain’s A record where the IP address is stored. It does this by locating the domain in the appropriate TLD nameserver based on the TLD attached to it. This is the .com TLD nameserver in the case of docs.google.com. Once it locates your second and top level domains, it looks for a subdomain as this may have a different IP address depending on how its DNS settings are configured. This means its search will trickle down to docs.google.com in the .com TLD nameserver for Google Docs. Once the DNS has found your record in the correct TLD nameserver, the authoritative server verifies the website’s identity via its IP address before returning it to the recursive resolver (from the original recursor server) so your browser can load the web page. You enter addresses into your browser regularly. Your browser does it for you when you use search engines and shortcuts. Either way, the DNS went through multiple steps to find the website’s exact location on the web for you.From your perspective, you seen a web page load within a few seconds in your browser. If you’ve already visited the website, the process is much shorter as the original recursive resolver will look through its cached information first to identify the website’s IP address rather than calling on the authoritative server. DNS servers explained The DNS recursor and authoritative servers often get confused with one another as they both return IP addresses to your browser. However, they’re quite different from one another. For instance, they’re used at different points in the DNS resolution process. The confusion stems from the recursor server’s ability to resolve DNS queries on its own. Normally, the recursor server acts as a liaison between your request and the authoritative server where the IP address is stored. However, when you’ve already visited a website and have not cleared your cache, the recursor server is able to return the site’s IP address on its own by reviewing its own cached data. Without that cached data, your query must travel down the DNS resolution pipeline like usual until it reaches the authoritative server. This server is the last step in the process as this server does not need to make additional requests. Its where DNS records are stored. If no record is found, it will return an error message instead, and you won’t be able to load the website you’re trying to visit. IP addresses are stored in different records within the authoritative server. You may have seen these records before if you’ve ever had to update DNS settings for your domain, such as when you want to connect an email client (like Google Workspace) to your domain. These records are comprised of multiple text files written in “DNS syntax.” Different records have different syntax, and each one has different instructions for how the authoritative server should handle the information contained within each record when requests come through. Here are the different types of records you’ll find attached to a domain and brief explanations for what they’re for: - A – Stores a domain’s IP address. - CNAME – Forwards an alias domain or subdomain to the actual domain it represents. CNAME records do not store IP addresses as they’re only used when the domains or subdomains stored within them are used as aliases for another domain. Alias domains do not have A records, so the authoritative server must forward requests to the A record of the domain the alias points to. - MX – Points to an email server. This is the record DNS servers use when you want to use your domain to send emails from business email addresses, such as email@example.com as opposed to firstname.lastname@example.org. - TXT – Used to store text notes from administrative purposes. - NS – The record used to store nameservers. This is what you’ll use when you want to register a domain with a dedicated registrar rather than your host. You’ll need to create a different NS record for each nameserver your host uses. The record points your domain to your host’s nameservers so the website you’ve stored there loads when you enter the domain attached to the record in a web browser. Many NS records also have “TTL” settings you can configure. This stands for “time to live,” or the number of times routers are able to pass the record around until it expires. It represents the number of times the recursor server can return a cached IP address it has stored. When the record expires (runs out of TTL counts), the server must send its request down the DNS resolution tube once more to find a domain’s IP address. You’ll also find TTL settings when you set up CDN caching. - SOA – Used to store admin information. TTL settings can be applied here as well. This record also contains information about admin email addresses and how long it’s been since the domain was updated. There are other DNS records, but these are the most common ones you’ll find attributed to your domain. The root nameserver is the first step in translating a domain name into its identifiable IP address. The recursor server sends its request here first. The root nameserver is responsible for passing that request onto the appropriate TLD nameserver. There are 13 types of root nameservers the DNS uses, and they’re all managed by a nonprofit organization called the Internet Corporation for Assigned Names and Numbers (ICANN). This organization controls all jurisdiction in regards to domains. It’s the organization that created the bylaw that requires you to attribute your personal information to every domain you register. Every recursive resolver is familiar with each type of root nameserver, and the DNS uses multiple copies of each around the world. Root nameservers are also responsible for applying Anycast routing to the traffic your domain receives when you use a CDN or registrar that supplies DDoS protection. Anycast is a network addressing method that routes traffic to multiple servers. This is as opposed to unicast routing, which sends traffic to a single server. TLD nameservers store information on domains based on the TLD each domain uses. For example, “docs.google.com” is stored in the .com TLD nameserver. Once the recursive resolver is sent to the correct TLD nameserver, it pinpoints the domain’s subdomain, if available, before the request is sent to the authoritative server. TLD nameservers are also overseen by ICANN, only these nameservers are managed by a branch of the organization called the Internet Assigned Numbers Authority (IANA). The IANA separates domains into two groups, gTLDs and ccTLDs, by combining gTLDs and sTLDs into one group. A lot of technical information is attributed to the domain name system. Fortunately, you won’t need to remember most of it in order to register and maintain your own domain. Still, you will need to update your domain’s nameservers if you don’t register it with your host. You’ll also need to update the DNS records if you want to use a CDN or business email clients. This doesn’t require much more than knowing where to copy and paste the correct records, which most services make easy through descriptive support tutorials. If you still need a domain, be sure to check out these guides:
ARTICLES IN THE BOOK This article is from: All text is available under the terms of the GNU Free Documentation License: http://en.wikipedia.org/wiki/Wikipedia:Text_of_the_GNU_Free_Documentation_License Dynamic range is a term used frequently in numerous fields to describe the ratio between the smallest and largest possible values of a changeable quantity. Dynamic range is an important indicator of the quality of a system intended either to record or to reproduce information for human perception. The human senses of sight and hearing have a very high dynamic range. A person is capable of hearing (and usefully discerning) anything from a quiet murmur in a soundproofed room to the melody in the loudest rock concert: the latter is 10,000,000,000 times louder than the former, that is a dynamic range of 100 dB. Equally a person can see objects in starlight (although colour differentiation is reduced at low light levels) or in bright sunlight, even though on a moonless night objects receive 1/1,000,000,000 of the illumination they would on a bright sunny day: that is a dynamic range of 90 dB. A person cannot perform these feats of perception at both extremes of the scale at the same time. The eyes take time to adjust to different light levels and the dynamic range of the human eye without any adjustment of the pupil is only approximately 30 dB. The instantaneous dynamic range of human audio perception is similar, so that, for example, a whisper cannot be heard in loud surroundings. Nevertheless, a good quality audio reproduction system should be able to reproduce accurately both the quiet sounds and the loud; and a good quality visual display system should be able to show both shadow details in nighttime scenes and the full brightness of sunny scenes. In practice it is difficult to achieve the full dynamic range seen by human beings using electronic equipment, since most electronic reproduction equipment is essentially linear rather than logarithmic like human perception. Electronically reproduced audio and video often uses some trickery to fit original material with a wide dynamic range into a narrower recorded dynamic range that can more easily be reproduced: this is dynamic compression. For example a good quality LCD display has a dynamic range of around 1000, or 30 dB (commercially the dynamic range is often called the "contrast ratio" meaning the full on/full off contrast ratio). When showing a movie or a game such a display is able to show both shadowy nighttime scenes and bright outdoor sunlit scenes, but in fact the level of light coming from the display is much the same for both types of scene (perhaps different by a factor of 10). Knowing that the display does not have a huge dynamic range, the program makers do not attempt to make the nighttime scenes millions of times less bright than the daytime scenes, but instead use other cues to suggest night or day: a nighttime scene will contain duller colours and will often be lit with blue lighting, which reflects the way that the human eye sees colours at low light levels. Audio engineers often use dynamic range to describe the ratio of the loudest possible undistorted sound to the quietest or to the noise level, say of a microphone or loudspeaker. In digital audio, the maximum possible dynamic range is given by the bit resolution (see signal-to-noise ratio). Dynamic range of an audio device is also sometimes referred to as the dynamic window. To mathematically determine a dynamic range you must add the headroom to the signal to noise ratio OR take the difference between the ceiling and noise floor of an audio device. For example, if the ceiling of a device is 10 dB and the floor is 3 dB then the dynamic range is 7 dB, since 10-3 = 7. Electronics engineers apply the term to: In audio and electronics applications, the ratio involved is often so huge that it is converted to a logarithm and specified in decibels. In music, dynamic range is the difference between the quietest and loudest volume of an instrument, part or piece of music. It is also the range of amplitudes an audio device can reproduce. Dynamic range is the headroom plus the signal to noise ratio ranges added together. It can be calculated by taking the difference between the ceiling and noise floor of an audio device. In modern recording, this range is often limited through audio level compression, which allows for louder volume, but can make the recording sound less exciting or live. Photographers use exposure range as a synonym for the luminosity range of a scene being photographed; the light sensitivity range of photographic film, paper and digital camera sensors; the opacity range of developed film images; the reflectance range of images on photographic papers. It can be controlled through the use of a graduated ND filter. More details about dynamic range and dynamic range optimization can be found here . In metrology, such as when performed in support of science, engineering or manufacturing objectives, “dynamic range” refers to the range of values that can be measured by a sensor or metrology instrument. Often this dynamic range of measurement is limited at one end of the range by saturation of a sensing signal sensor or by physical limits that exist on the motion or other response capability of a mechanical indicator. The other end of the dynamic range of measurement is often limited by one or more sources of random noise or uncertainty in signal levels that may be described as the defining the sensitivity of the sensor or metrology device. When digital sensors or sensor signal converters are a component of the sensor or metrology device, the dynamic range of measurement will be also related to the number of binary digits (“bits”) into which any analog measurement quantities are converted to create digital numeric values. For example, a 12-bit digital sensor or converter can only provide a dynamic range in which the ratio of the maximum measured value to the minimum measured value is limited to 4096-to-1. Metrology systems and devices may use several basic methods to increase their basic dynamic range. These methods include averaging and other forms of filtering, repetition of measurements, nonlinear transformations to avoid saturation, etc. In more advance forms of metrology, such as multiwavelength digital holography, interferometry measurements made at different scales (different wavelengths) can be combined to retain the same low-end resolution while extending the upper end of the dynamic range of measurement by orders of magnitude. High Dynamic Range Imaging is an emerging field in computer graphics which seeks to represent light levels (either measured or synthesised) as an open-ended range of absolute values, rather than as a simple ratio of 'full' brightness. This allows more accurate and realistic renderings. Standard Operating Level: A specified reference level. In recording applications, standard operating level is defined as O VU = + 4 dBm.
Hearing loss is very common. It can be very distressing, particularly if it is getting worse and especially if it affects both ears. Everyone's hearing gets worse as they get older. However, there are many types of hearing loss and not all are restricted to older people. In this article we talk about the different causes of hearing loss, their different characteristics and what you can do about them. Structure of the ear The ear is roughly divided into three parts. The outer (external) ear includes the part you can see (called the pinna) and the narrow tube-like structure (the ear canal), which your health professional can look down with a torch. At the end of the canal is the eardrum. This separates the external ear from the middle ear. The eardrum is a tightly stretched membrane, a bit like the skin of a drum. The middle ear is an air-filled compartment. Inside it are the three smallest bones in the body, called the malleus, incus and stapes. These bones are connected to each other. The last in the line, the stapes, also makes contact with the inner (internal) ear. The air space of the middle ear connects to the back of the nose by the Eustachian tube. The inner ear is made up of two components, the cochlea and the vestibular system. The cochlea is involved with hearing. The vestibular system helps with balance. The cochlea is a snail-shaped chamber filled with fluid. It is lined with special sensory cells called hair cells. These cells transform sound waves into electrical signals. The cochlea is attached to a nerve that leads to the brain. The vestibular system is made up of a network of tubes, called the semicircular canals, plus the vestibule. The vestibular system also contains special sensory cells but here they detect movement instead of sound. Both the cochlea and the vestibular system are connected to a nerve which carries electrical signals to the brain. How do we hear? Sound waves are created when air vibrates. To hear, the ear must change sound into electrical signals which the brain can interpret. The outer part of the ear (the pinna) funnels sound waves into the ear canal. When sound waves reach the eardrum they make it vibrate. Vibrations of the eardrum make the tiny bones in the middle ear move too. The last of these bones (the stapes) passes on the vibrations to the fluid-filled chamber called the cochlea. When the vibrations reach the cochlea, the fluid inside it moves. As the fluid moves it vibrates the hairs on the cells that line the cochlea. Each cell is stimulated by a particular note (or frequency) of sound. The vibration of the hair cells is turned into an electrical signal by the organ of Corti, at their base. The organ of Corti then sends signals down the hearing (auditory) nerve to the brain. Special areas in the brain receive these signals and translate them into what we know as sound. Your ears create electrical signals that represent an extraordinary variety of sounds. For example, the speed at which the eardrum vibrates varies with different types of sound. With low-pitched sounds the eardrum vibrates slowly. With high-pitched sounds it vibrates faster. This means that the special hair cells in the cochlea also vibrate at varying speeds. This causes different signals to be sent to the brain. This is one of the ways in which we are able to distinguish between a wide range of sounds. What can cause hearing loss? Damage to any part of the ear can cause a hearing loss. Conductive hearing loss If there is a problem in the ear canal or the middle ear, this causes what is known as a conductive hearing loss. In conductive hearing loss, the movement of sound (conduction) is blocked or does not pass into the inner ear. This is often as the result of earwax (cerumen) or fluid in the middle ear, although it may also be caused by a burst (ruptured) eardrum or by otosclerosis (see below). Sensorineural hearing loss If the fluid-filled chamber called the cochlea or the hearing nerve is not working properly this causes what is known as a sensorineural hearing loss. Usually this means that hair cells in the cochlea are not working properly or there is a problem with the hearing nerve so that some or all sounds are not being sent to the brain. It does not usually affect the entire range of sound frequencies, at least not at first. Sensorineural hearing losses are usually permanent. They can be mild, moderate, severe or profound and affect one or both ears. It is also possible for sensorineural and conductive hearing losses to occur together in a mixed hearing loss. Causes of conductive hearing loss: the eardrum and ear canal Blockage of the ear canal The most common cause of blockage in the ear canal is earwax (cerumen). See the separate leaflet called Earwax for more details. Something in the ear canal Objects that shouldn't be there (foreign bodies) are most often found in the ears of children. Peas, beads or small pieces of a toy are the most common foreign bodies to block the ear and affect hearing. It is usually best to have the object removed. This is done either by syringing it out with warm water or with a special extracting device. Sometimes the skin of the ear canal can become inflamed. This may be caused by infection, allergy or other causes. Common symptoms include itch, ear discharge and dulled hearing. It is treated with ear drops. See the separate leaflet called Ear Infection (Otitis Externa) for more details. A torn (perforated) eardrum is not usually serious and often heals on its own without any complications. It may cause hearing loss, in which case a small procedure to repair it is an option. See the separate leaflet called Perforated Eardrum for more details. Scarring of the drum is usually due to repeated perforation, either by infection or by poking objects into the ear. Cholesteatoma is an uncommon condition where a growth develops in the ear. You can be born with it but usually it occurs as a complication of a long-standing (chronic) ear infection. The most common symptoms are a loss of hearing and a smelly discharge from the ear. See the separate leaflet called Cholesteatoma for more details. The need for an eardrum Not having an eardrum does not result in deafness. Part of the job of the eardrum is to boost (amplify) sound. Without the eardrum the sound will still reach the middle ear; however, it will not be as loud. Its other job is to seal off the middle ear and prevent it from damage by water and soap. Causes of conductive hearing loss: the middle ear The middle ear consists of an air space and the three small hearing bones (ossicles). Conduction of sound through the middle ear depends upon both of these. The air in the middle ear gets there from the Eustachian tube, which links the middle ear to the throat. The tube allows air in and out (the ear popping sensation on a flight shows this happening). This allows the air pressure inside the ear to equalise with that outside and replaces air that gets absorbed by the cells lining the middle ear. The middle ear system can therefore be affected by problems affecting the Eustachian tubes and the middle ear space itself and by problems affecting the hearing bones. Conditions affecting the air space Eustachian tube dysfunction When your Eustachian tube isn't working properly this can dull your hearing. It is usually a temporary problem that lasts a week or so and most commonly happens during or after a cold. The middle ear can fill with fluid. This is called a middle ear effusion. See the separate leaflet called Eustachian Tube Dysfunction for more details. Glue ear happens when the middle ear fills with a glue-like liquid instead of air. Usually it clears without treatment. However, an operation to clear the fluid and put in a tiny tube (grommet) to allow air to get into the middle ear may be advised if it persists. See the separate leaflet called Glue Ear for more details. An ear infection is very common, particularly in children. The main symptoms are earache and feeling unwell but it can cause temporary hearing loss. See the separate leaflet called Ear Infection (Otitis Media) for more details. Conditions affecting the hearing bones Otosclerosis is the most common cause of hearing loss in young people. It mainly affects the third of the three bones in the middle ear (the stapes). It causes gradual hearing loss. Treatments include hearing aids and surgery. See the separate leaflet called Otosclerosis for more details Causes of sensorineural hearing loss: the cochlea Hearing loss of older people (presbyacusis) The most common cause of hearing loss is age-related. Most people over the age of 60 develop hearing loss to some degree. The exact cause is not known but it is thought to be due to the cells in the cochlea becoming damaged over time. A hearing aid may be needed. See separate leaflet called Hearing Loss of Older People (Presbyacusis) for more details. Loud noise damages the cochlea and can result in permanent hearing loss and ringing in the ears (tinnitus). The risk is based on how loud the noise is and how long you have been exposed to it. Those who work with loud equipment - people who shoot, use pneumatic drills or operate heavy machinery - should always wear their protective ear-wear in order to prevent long-term damage. The cumulative effect of prolonged exposure to loud noise speeds up the process of hearing loss. Loud noise from MP3 players and music gigs is thought to be the reason why hearing loss is increasingly affecting young people. If you ever have ringing in your ears or dull hearing after listening to music, it was too loud; many musicians now wear ear filters to protect their hearing. Other causes of cochlear damage The cochlea can be damaged by a severe head injury. Such trauma can also disrupt the tiny ear bones (the ossicles) and cause hearing loss that way. The cochlea can be damaged by a cholesteatoma (see above). Some infections may damage the hearing nerves and/or the cochlea. These include measles, mumps, bacterial and viral meningitis and tuberculosis. The zoster virus, which causes shingles, can affect the hearing nerves. In all of these cases the resulting loss of hearing can be permanent. The cochlea is also vulnerable to damage by poisonous substances (toxins). This includes some medicines which can damage the hearing. Certain antibiotics are known to carry a risk of this but they may still be used where nothing else will work, particularly if life is at risk. The small risk of hearing nerve damage is felt to be outweighed by the need for the medicine. One example is gentamicin, an antibiotic which is invaluable in severe infections by particular germs (bacteria), especially in babies. Any baby who is treated with gentamicin will automatically have a hearing nerve testing once they have recovered from the infection. Ménière's disease causes attacks of dizziness and tinnitus, as well as hearing loss. See separate leaflet called Ménière's Disease for more details. Causes of sensorineural hearing loss: the auditory nerve and the brain Conditions affecting the hearing (auditory) nerve An acoustic neuroma is a rare growth on the hearing nerve inside the skull. The hearing loss that an acoustic neuroma causes affects just one ear. See separate leaflet called Acoustic Neuroma for more details. Conditions affecting the brain Ultimately, sound is heard and interpreted by the brain; conditions that affect the hearing centre in the brain can also cause hearing loss. Examples would be brain injury through trauma, stroke, brain infection (encephalitis) and multiple sclerosis. In some cases there might be partial or even complete recovery over time, although in other cases the loss would be permanent. Congenital hearing loss Congenital hearing loss is hearing loss present at birth or soon after. Most of it is inherited but about 1 in 4 cases are due to things that happen to the baby in the womb. This includes infections (such as German measles or cytomegalovirus) and being premature or not getting enough oxygen at birth. Microtia is a condition in which the ear is underdeveloped or does't develop at all. There are many inherited (genetic) syndromes which may cause hearing loss. These include Down's syndrome, Treacher Collins' syndrome and Waardenburg's syndrome. They may have their effects through the mechanisms described in this leaflet, although in some cases they are due to abnormal development of the hearing apparatus in the womb. Some inherited hearing loss in children is not present at birth but develops in the early years. The first sign can be poor speech development. In the UK all newborn babies have a hearing test to make sure the hearing (auditory) nerve is working properly. Babies are checked again at about 8 months of age to make sure they can still hear. Families are encouraged to talk to health visitors or doctors if their child's speech seems delayed; most cases are detected fairly quickly. In many cases treatment to improve hearing is possible. This might involve hearing aids, cochlear implants or tiny tubes called grommets. In cases of profound deafness, where hearing cannot be improved, children and their families are taught how to use sign language from a very early age in order to help communication. Unilateral versus bilateral hearing loss: one ear or two? The most common causes of hearing loss, such as age-related hearing loss, tend to affect both ears. However, this is often unequal, with the hearing loss of the ears differing from one another. Some people will have hearing aids fitted to just one ear, whereas others will choose to aid both ears. Loss of hearing in just one ear is sometimes called single-sided deafness (SSD) or hearing loss. Being deaf in one ear presents particular challenges: - It affects directional hearing - that is, the ability to work out which direction sound is coming from. This is also called spatial hearing. If you have normal hearing, you use the time difference in a sound arriving at each ear and the difference in loudness between two ears, to work out where a sound is coming from and how far away it is. If you have significant hearing loss in just one ear this makes it much harder to localise a sound. For example, hearing someone call you outside or hearing whether a car is coming when you are about to cross the road. - When a sound is coming from your affected side, your head gets in the way of the sound getting to your good ear - the head shadow effect. This is particularly noticeable with higher-pitched sounds such as the c, f, p, s, t, ch and sh sounds, making it harder to discriminate between some words. This is why people with single-sided hearing loss have problems hearing speech when there is background noise, even when the hearing is normal in their other ear. If you are affected by one-sided hearing loss, here are some coping strategies that may help: - Don't be afraid to tell people that you are deaf in one ear or that you hear much better in one ear than in the other. - Always make sure people are seated or walk on your good side. - If there is background noise coming from only one direction, put yourself so that the noise is to your bad ear. - If your affected ear is on the passenger side of your car, seating someone behind you rather than next to you may help. - If you are going to a meeting, arrive early enough to choose a seat so that most people are on your good side. (This is easier at a rectangular rather than a round table - for example, in the corner with your back to a wall.) - Remind your friends, family, teachers and co-workers not to expect you to answer if they call your name in a crowded place or across a road; otherwise they may think you are brusque or aloof and ignoring them. - Keep your mobile phone in your pocket, as locating it when it rings can be very hard. - Use a mono-splitter to listen to stereo music. - Enjoy being able to sleep in a noisy environment by sleeping on your 'good' side! Bone-anchored hearing aid (BAHA) and contralateral routing of signal (CROS) hearing aids are specifically designed for people with one-sided hearing loss. They pipe the sounds from your 'bad ear' side to your good or better side. Protecting your hearing The best way to protect your hearing is to avoid exposure to very loud noise completely and to limit the amount of time you are exposed to loud noise. Hearing loss can be disabling and can lead to feelings of isolation and depression. Discussing your hearing loss with your doctor If you think you are losing hearing then discuss this with your GP. You may notice difficulty distinguishing what people are saying, or that everything seems quieter. Others may comment that you have the TV turned up very loud. Perhaps you have been exposed to loud noise. Your doctor may ask you the following questions: - Has your hearing worsened suddenly or gradually? Sudden deafness is uncommon but can be due to sudden trapping of water behind wax (cerumen) in the ear canals. Rarely it can have a more serious cause and you may need to see a hearing specialist urgently. - Are both ears affected? Clearly ear infections, glue ear and earwax may be one-sided or two-sided. However, if you have severe hearing loss and wax is present in the ear canals then the wax needs to be cleared. After this the ear needs to be reassessed in case the wax was not the only reason for hearing reduction. - Do you have associated sinus problems/congestion? If so glue ear is very likely. This may be following a cold, when mucus remains in the middle ear space for several weeks. It may also be associated with allergic conditions such as hay fever. - Have you been exposed to loud noise? Are you aged over 50 years? Either of these increases the chance that this is age-related high-frequency hearing loss. - Do you have particular difficulties with hearing conversation? Again this is a feature of high-frequency hearing loss, as consonants, which make language clear, are higher-frequency sounds than vowels. - Is there a family history of hearing loss? Some causes of deafness run in families, particularly otosclerosis. Your doctor is likely to perform some tests on your hearing in the surgery to decide whether your hearing loss seems to be conductive, sensorineural, or mixed. They will examine your ears for wax and for obvious problems affecting the eardrum. They may then refer you for formal hearing tests. See the separate leaflet called Hearing Tests for more details. There are many causes of hearing loss: the most common in younger people is otosclerosis, whereas the most common in older people is hearing loss of older people (presbyacusis). Many people have mixed hearing loss: you may have hearing loss of older people and earwax (cerumen) for instance, or otosclerosis and earwax. It's very important, if your doctor finds a simple cause for your reduced hearing and treats it, that you return to your doctor if your hearing is still not better. Many conditions are treatable and a very few need urgent management. Further reading and references ; NICE Guideline (July 2018) ; Hearing aids and how to get one, ENT UK, 2011 ; Presbycusis. Lancet. 2005 Sep 24-30366(9491):1111-20. ; Noise Help So recently about 2 weeks ago im getting this constant ringing or buzzing noise in my ear. Before it use to just come once in awhile now its almost an all day thing. My family doctor says looks like...brandon09344 Disclaimer: This article is for information only and should not be used for the diagnosis or treatment of medical conditions. Patient Platform Limited has used all reasonable care in compiling the information but make no warranty as to its accuracy. Consult a doctor or other health care professional for diagnosis and treatment of medical conditions. For details see our conditions.
A high-speed centrifuge is a type of separator most often used on marine vessels and power generation stations to remove contamination in fuel and lubricating oils such as solids and water. It is imperative to carry out this treatment to prevent damage to engines and generators. The separation principle of high-speed centrifuge works on the differing specific gravity between two different liquids. To understand, we examine a beaker with oil, water and solids. After a settling period, water (heavy liquid) and solids get collected at the bottom while the oil (light liquid, non-emulsified) floats up top. take note that the separation is due to gravity (or specific gravity). Heavy bunker fuel oil has an SG of about 0.95, diesel oil about 0.85 and fresh water has an SG of 1. Because of the difference in Specific Gravity, or SG, the oil will float on top of the water. The solid particles that is heavier than water will sink down. But by gravity, separation will take a long time. Furthermore, if the SG's of the mixture are very close, the oil and the water may not be able to separate very well. Mathematically this process can be represented by: Fs = ∏/6x D3 (ρw-ρo) g Where Fs is the separating force, ρw is the density of water, ρo is the density of oil and “g” is gravitational force. If we position the beaker on its side and rotate it fast, then the gravitational factor g will be replaced by the centrifugal force ω2 r, where ω2 is angular velocity of rotation and r is the effective radius. Fs = ∏/6x D3 (ρw-ρo) ω2 r. Now the separating force will be much higher in the centrifuge as compare to a beaker. The separation effect is achieved much faster with a much greater force. As illustrated in the figure by positioning the beaker on its side and rotating it about the top as its axis. This movement forms the basic principle of centrifugal purifiers. Through a system of gears, a centrifuge bowl is rotated to and maintained at high speeds. Oil to be cleaned is allowed to enter the bowl while it is rotating. The heavier components in the oil are thus forced outwards. The solid particles that are too fine to be removed by filtration are forced towards the circumference of the bowl. Below is an illustration of the process during operation of the Mitsubishi Selfjector Hercules Series oil separators. The oil should also be heated pre-feed to lower the SG of the oil. The difference in SG's between the oil and the water will thus be wider; making for better separation between the oil and the water. Oil purifiers usually maintain a layer of water inside the bowl to act as a seal for the oil known as sealing water. Without it, the oil will flow out together with the contaminants. If removal of water is not needed, the centrifuge can be modified so that no water layer is needed. The centrifuge then becomes a clarifier. Industmarine Engineers’ continual distinction in embodying our core philosophies of innovation, the human touch and active engagement has helped establish our company as a regional leader for providing products and services pertaining to the marine industry.
Table of Contents In this article, we will explore how to calculate the sum of array elements using recursion in C++. We will provide a clear explanation of the concept and then walk through a step-by-step example. Recursion is a powerful programming technique that allows a function to call itself in order to solve a problem. It is often used to simplify complex problems and can be a valuable tool in a programmer’s toolkit. Recursion involves breaking down a problem into smaller, more manageable subproblems. In the context of summing array elements, we can think of the problem as follows: - If the array is empty (i.e., it has no elements), the sum is 0. - If the array has one element, the sum is that element. - If the array has more than one element, we can break it down into two parts: - The first element (the head) and - The remaining elements (the tail). We can calculate the sum of the array by adding the head to the sum of the tail, where the sum of the tail is calculated using the same recursive function. Sum of Array Elements Using Recursion Explanation of the Code - We define a recursive function sumArraythat takes two arguments: an array arrand its size - In the base case, when sizeis 0 (i.e., the array is empty), the function returns 0 because there are no elements to sum. - In the recursive case, we calculate the sum of the array by adding the first element arrto the result of the recursive call sumArray(arr + 1, size - 1), where arr + 1represents the tail of the array, and size - 1indicates the reduced size of the array. - In the mainfunction, we create an example integer array arrand calculate its sum using the sumArrayfunction. We then print the result to the console.
Inversion (discrete mathematics) The inversion is usually defined for permutations, but may also be defined for sequences: Let be a sequence (or multiset permutation). If and , either the pair of places or the pair of elements is called an inversion of . For sequences, inversions according to the element-based definition are not unique, because different pairs of places may have the same pair of values. The inversion set is the set of all inversions. A permutation's inversion set according to the place-based definition is that of the inverse permutation's inversion set according to the element-based definition, and vice versa, just with the elements of the pairs exchanged. The inversion number of a sequence , is the cardinality of inversion set. It is a common measures of (pre-)sortedness of a permutation or sequence. This value is comprised between 0 and included. For example since the sequence is ordered. Also as each pairs is an inversion. This last example shows that a sort that is intuitively sorted can still have a quadratic number of inversions. It does not matter if the place-based or the element-based definition of inversion is used to define the inversion number, because a permutation and its inverse have the same number of inversions. Other measures of (pre-)sortedness include the minimum number of elements that can be deleted from the sequence to yield a fully sorted sequence, the number and lengths of sorted "runs" within the sequence, the Spearman footrule (sum of distances of each element from its sorted position), and the smallest number of exchanges needed to sort the sequence. Standard comparison sorting algorithms can be adapted to compute the inversion number in time O(n log n). Three similar vectors are in use that condense the inversions of a permutation into a vector that uniquely determines it. They are often called inversion vector or Lehmer code. (A list of sources is found here.) This article uses the term inversion vector () like Wolfram. The remaining two vectors are sometimes called left and right inversion vector, but to avoid confusion with the inversion vector this article calls them left inversion count () and right inversion count (). Interpreted as a factorial number the left inversion count gives the permutations reverse colexicographic, and the right inversion count gives the lexicographic index. Inversion vector : With the element-based definition is the number of inversions whose smaller (right) component is . - is the number of elements in greater than before . Left inversion count : With the place-based definition is the number of inversions whose bigger (right) component is . - is the number of elements in greater than before . Right inversion count , often called Lehmer code: With the place-based definition is the number of inversions whose smaller (left) component is . - is the number of elements in smaller than after . Both and can be found with the help of a Rothe diagram, which is a permutation matrix with the 1s represented by dots, and an inversion (often represented by a cross) in every position that has a dot to the right and below it. is the sum of inversions in row of the Rothe diagram, while is the sum of inversions in column . The permutation matrix of the inverse is the transpose, therefore of a permutation is of its inverse, and vice versa. Example: All permutations of four elements The following sortable table shows the 24 permutations of four elements with their place-based inversion sets, inversion related vectors and inversion numbers. (The small columns are reflections of the columns next to them, and can be used to sort them in colexicographic order.) It can be seen that and always have the same digits, and that and are both related to the place-based inversion set. The nontrivial elements of are the sums of the descending diagonals of the shown triangle, and those of are the sums of the ascending diagonals. (Pairs in descending diagonals have the right components 2, 3, 4 in common, while pairs in ascending diagonals have the left components 1, 2, 3 in common.) The default order of the table is reverse colex order by , which is the same as colex order by . Lex order by is the same as lex order by . Weak order of permutations If a permutation is assigned to each inversion set using the place-based definition, the resulting order of permutations is that of the permutohedron, where an edge corresponds to the swapping of two elements with consecutive values. This is the weak order of permutations. The identity is its minimum, and the permutation formed by reversing the identity is its maximum. If a permutation were assigned to each inversion set using the element-based definition, the resulting order of permutations would be that of a Cayley graph, where an edge corresponds to the swapping of two elements on consecutive places. This Cayley graph of the symmetric group is similar to its permutohedron, but with each permutation replaced by its inverse. \|Wikiversity has learning resources about Inversion (discrete mathematics)\| \|Wikimedia Commons has media related to Inversion (discrete mathematics).\| - Factorial number system - Permutation graph - Transpositions, simple transpositions, inversions and sorting - Damerau–Levenshtein distance - Parity of a permutation Sequences in the OEIS: - Sequences related to factorial base representation - Factorial numbers: A007623 and A108731 - Inversion numbers: A034968 - Inversion sets of finite permutations interpreted as binary numbers: A211362 (related permutation: A211363) - Finite permutations that have only 0s and 1s in their inversion vectors: A059590 (their inversion sets: A211364) - Number of permutations of n elements with k inversions; Mahonian numbers: A008302 (their row maxima; Kendall-Mann numbers: A000140) - Number of connected labeled graphs with n edges and n nodes: A057500 - Aigner 2007, pp. 27. - Comtet 1974, pp. 237. - Knuth 1973, pp. 11. - Pemmaraju & Skiena 2003, pp. 69. - Vitter & Flajolet 1990, pp. 459. - Bóna 2012, pp. 57. - Cormen et al. 2001, pp. 39. - Barth & Mutzel 2004, pp. 183. - Gratzer 2016, pp. 221. - Mannila, 1984 & pp318. sfn error: no target: CITEREFMannila1984pp318 (help) - Mahmoud 2000, pp. 284. - Kleinberg & Tardos 2005, pp. 225. - Weisstein, Eric W. "Inversion Vector" From MathWorld--A Wolfram Web Resource - Reverse colex order of finitary permutations (sequence A055089 in the OEIS) - Aigner, Martin (2007). "Word Representation". A course in enumeration. Berlin, New York: Springer. ISBN 978-3642072536. - Barth, Wilhelm; Mutzel, Petra (2004). "Simple and Efficient Bilayer Cross Counting". Journal of Graph Algorithms and Applications. 8 (2): 179–194. doi:10.7155/jgaa.00088. - Bóna, Miklós (2012). "2.2 Inversions in Permutations of Multisets". Combinatorics of permutations. Boca Raton, FL: CRC Press. ISBN 978-1439850510. - Comtet, Louis (1974). "6.4 Inversions of a permutation of [n]". Advanced combinatorics; the art of finite and infinite expansions. Dordrecht,Boston: D. Reidel Pub. Co. ISBN 9027704414. - Cormen, Thomas H.; Leiserson, Charles E.; Rivest, Ronald L.; Stein, Clifford (2001). Introduction to Algorithms (2nd ed.). MIT Press and McGraw-Hill. ISBN 0-262-53196-8. - Gratzer, George (2016). "7-2 Basic objects". Lattice theory. special topics and applications. Cham, Switzerland: Birkhäuser. ISBN 978-3319442358. - Kleinberg, Jon; Tardos, Éva (2005). Algorithm Design. ISBN 0-321-29535-8. - Knuth, Donald (1973). "5.1.1 Inversions". The art of computer programming. Addison-Wesley Pub. Co. ISBN 0201896850. - Mahmoud, Hosam Mahmoud (2000). "Sorting Nonrandom Data". Sorting: a distribution theory. Wiley-Interscience series in discrete mathematics and optimization. 54. Wiley-IEEE. ISBN 978-0-471-32710-3. - Pemmaraju, Sriram V.; Skiena, Steven S. (2003). "Permutations and combinations". Computational discrete mathematics: combinatorics and graph theory with Mathematica. Cambridge University Press. ISBN 978-0-521-80686-2. - Vitter, J.S.; Flajolet, Ph. (1990). "Average-Case Analysis of Algorithms and Data Structures". In van Leeuwen, Jan (ed.). Algorithms and Complexity. 1 (2nd ed.). Elsevier. ISBN 978-0-444-88071-0. - Margolius, Barbara H. (2001). "Permutations with Inversions". Journal of Integer Sequences. 4. - Mannila, Heikki (1984). "Measures of presortedness and optimal sorting algorithms". Lecture Notes in Computer Science. 172: 324–336. doi:10.1007/3-540-13345-3_29. ISBN 978-3-540-13345-2. - Estivill-Castro, Vladimir; Wood, Derick (1989). "A new measure of presortedness". Information and Computation. 83 (1): 111–119. doi:10.1016/0890-5401(89)90050-3. - Skiena, Steven S. (1988). "Encroaching lists as a measure of presortedness". BIT. 28 (4): 755–784. doi:10.1007/bf01954897. S2CID 33967672.
An algorithm includes finitely many steps to solve a given task. If the sequence of these operations is followed, the task will be solved subsequently in any case and the same results will always be delivered. Algorithms are widely used, especially in mathematics and computer science, because they can be used to solve complex problems easily and repeatedly. How is the term defined? We encounter algorithms in more and more areas of everyday life. Although the term is primarily at home in computer science and mathematics, we also encounter the simplest algorithms in everyday life. The hand-washing instructions which have been posted on many public sinks since the Corona pandemic are also an example. For everyday use, the algorithm can therefore be understood as a fixed sequence of steps or work instructions, the execution of which leads to the desired result, namely in this case clean and hygienic hands. In even more general terms, they define procedures by which input values (dirty hands) can be transformed into fixed output values (clean/hygienic hands). In computer science, however, this general definition is not sufficient. An algorithm is a finite set of steps whose execution leads to the solution of a predefined problem. If it is executed several times, it always delivers the same result. Furthermore, an algorithm is well-defined if it contains only unambiguous instructions that can be executed exactly as they are. For example, a step of the hand washing instructions could be “Wash your hands under the stream of water.” This step is not well-defined because the execution is not clear from it. For example, questions such as “For how long do the hands need to be washed?” or “How strong should the water jet be for this?” still arise. What are the Properties of an Algorithm? The following properties are indicative: - Unambiguity: The descriptions and steps must be unambiguous. - Executability: All steps must be executable if the previous steps were executed correctly. - Finiteness: Finiteness means that there is a finite set of steps to be executed. - Termination: Termination means that the algorithm reaches a result after a certain amount of time. Termination and finiteness are mutually dependent. So by the finiteness of the steps, the function must come also forced to an end. - Determinacy: Given the same circumstances, the process steps always lead to the same result. - Determinism: At each time of the execution there is exactly one unique subsequent step, which must be executed. This property is also a prerequisite for determinacy. If there were a selection of possible subsequent steps, there could not be an unambiguous result. What are some examples? - Navigation Device: The route description of a navigation device contains many successive steps that eventually lead to the entered destination. It thus satisfies all the properties of algorithms. - Game Roboters: When computers learn games, such as chess or Go, they in many cases adhere to predefined sequences of steps. These have been programmed so that they are simply executed when a certain situation is recognized. - Mathematics: Mathematical functions also represent predefined sequences of steps, since they specify a precise sequence of how to get from an input value to an output value. - Traffic Lights: The switching times of traffic lights are also determined by clear instructions. Depending on the time or the measured traffic density, the traffic light switches from red to green and vice versa. What are the common algorithms used in computer science? Common algorithms are used in practice in a variety of fields such as computer science, engineering, statistics, etc. Some of the most commonly used are: - Sorting: Sorting processes are used to arrange a list of data in a particular order. Some of the most common sorting algorithms include bubble sort, insertion sort, selection sort, merge sort, and quicksort. - Search: Search processes are used to find a specific value or element in a list of data. The most commonly used ones include linear search, binary search, and interpolation search. - Graph: Graph algorithms are used to analyze and process graphs, which are structures consisting of nodes and edges. The most commonly used graph algorithms include the one defined by Dijkstra for finding the shortest paths, the one by Kruskal for minimum spanning trees, and Floyd-Warshall’s algorithm for shortest paths for all pairs. - Machine Learning: Machine learning algorithms are used in artificial intelligence to automatically improve performance based on data. The most commonly used ones include linear regression, logistic regression, decision trees, and neural networks. - Compression: These are used to reduce the size of data for efficient storage and transmission. The most commonly used include Huffman coding, Lempel-Ziv-Welch (LZW) algorithm, and run-length coding. - Encryption: Encryption processes are used to secure data and protect it from unauthorized access. The most commonly used encryption algorithms include Advanced Encryption Standard (AES), Rivest-Shamir-Adleman (RSA) encryption, and Blowfish. - Hashing: Hashing algorithms are used to map data of arbitrary size to fixed-size values. The most commonly used ones include MD5, SHA-1, and SHA-256. These and many other algorithms are used in various applications such as web search engines, image and video processing, data analysis, networking, security, etc. If programmers and data scientists understand the basics of these, they can design and implement efficient and effective software solutions. How are Machine Learning and Algorithms related? As we have already learned, these are defined process chains that always achieve the same result with the same inputs. In the field of Machine Learning, it is therefore often referred to as algorithms that are capable of learning complex relationships. For the grouping of data points, the so-called k-means clustering is often used. Here, different cluster centers are tested in a finite set of steps until sooner or later the perfect assignment in groups is learned. So this is actually an algorithm. In the field of deep learning, however, this assignment is not quite so obvious. In the case of neural networks, two training runs with the same data sets can lead to different results. This would actually speak against an assignment to the algorithms. However, this is mainly due to the initialization of the network during which, for example, the weights of the individual neurons are randomly assigned. If among other things, it is ensured that the weights are the same in two training runs, the results of the networks would also be the same. Thus, deep learning models also count as algorithms. Are algorithms fair? Algorithmic bias and fairness have become increasingly important topics in recent years as more and more decisions are made by them. The bias occurs when an algorithm produces results that systematically and unfairly discriminate against certain groups of people based on race, gender, age, or other characteristics. One of the biggest concerns is that it can perpetuate and even exacerbate existing social inequalities. For example, if such a system is used in hiring decisions and is biased against women or minorities, it can lead to a perpetuation of a lack of diversity in the workplace. To address it, it is important to understand how it can occur. One common cause of bias is biased training data. If the data used to train an algorithm is biased, it will reflect that bias in its results. This is especially true of Machine Learning algorithms, which learn from historical data. Another cause of bias is the assumptions made by the developers. If the developers have implicit biases, those biases can be reflected in the process steps. To mitigate the bias, several approaches can be taken. One approach is to use diverse training data that accurately reflects the diversity of the population. Another approach is to design code that is explicitly fair, such as those that use statistical methods to ensure that decisions are not based on protected characteristics. It is important to note that while algorithms can be designed to be fair, there is no single definition of fairness. Different groups may have different views on what constitutes fair decision-making, and it is important to consider those perspectives when designing these systems. In conclusion, algorithmic bias and fairness are critical considerations that impact people’s lives. By being aware of the potential for bias and taking steps to mitigate it, we can create systems that are more equitable and just. This is what you should take with you - An algorithm includes finitely many steps to solve a given task. - If the sequence of these workflows is followed, the task will be solved subsequently in any case and the same results will always be delivered. - We encounter algorithms very often in everyday life, such as in navigation devices or traffic lights. Machine learning models are also among them. Other Articles on the Topic of Algorithms On Wikipedia, there is a detailed article on the topic.
In computer programming, an integer overflow occurs when an arithmetic operation attempts to create a numeric value that is outside of the range that can be represented with a given number of bits - either larger than the maximum or lower than the minimum representable value. The most common result of an overflow is that the least significant representable bits of the result are stored the result is said to wrap around the maximum (i.e. modulo power of two). An overflow condition gives incorrect results and, particularly if the possibility has not been anticipated, can compromise a program's reliability and security. On some processors like graphics processing units (GPUs) and digital signal processors (DSPs), the result saturates; that is, once the maximum value is reached, any attempt to increase it always returns the maximum integer value. - 8 bits: maximum representable value 28 − 1 = 255 - 16 bits: maximum representable value 216 − 1 = 65,535 - 32 bits: maximum representable value 232 − 1 = 4,294,967,295 (the most common width for personal computers as of 2005[update]), - 64 bits: maximum representable value 264 − 1 = 18,446,744,073,709,551,615 (the most common width for personal computer CPUs, as of 2017[update]), - 128 bits: maximum representable value 2128 − 1 = 340,282,366,920,938,463,463,374,607,431,768,211,455 When an arithmetic operation produces a result larger than the maximum above for a N-bit integer, an overflow reduces the result to modulo N-th power of 2, retaining only the least significant bits of the result and effectively causing a wrap around. In particular, multiplying or adding two integers may result in a value that is unexpectedly small, and subtracting from a small integer may cause a wrap to a large positive value (for example, 8-bit integer addition 255 + 2 results in 1, which is 257 mod 28, and similarly subtraction 0 - 1 results in 255, a two's complement representation of -1). Such wrap around may cause security problems - if an overflown value is used as the number of bytes to allocate for a buffer, the buffer will be allocated unexpectedly small, leading to a potential buffer overflow and arbitrary code execution. If the variable has a signed integer type, a program may make the assumption that a variable always contains a positive value. An integer overflow can cause the value to wrap and become negative, which violates the program's assumption and may lead to unexpected behavior (for example, 8-bit integer addition of 127 + 1 results in -128, a two's complement of 128). Most computers have two dedicated processor flags to check for overflow conditions. The carry flag is set when the result of an addition or subtraction, considering the operands and result as unsigned numbers, does not fit in the given number of bits. This indicates an overflow with a carry/borrow from the most significant bit. An immediately following add with carry or substract with borrow operation would use the contents of this flag to modify a register or a memory location that contains the higher part of a multi-word value. The overflow flag is set when the result of an operation on signed numbers does not have the sign that one would predict from the signs of the operands, e.g. a negative result when adding two positive numbers. This indicates than an overflow has occurred and the signed result represented in two's complement form would not fit in the given number of bits. Methods to mitigate integer overflow problems \|Language\|\|Unsigned integer\|\|Signed integer\| \|Ada\|\|modulo the type’s modulus\|\|raise Constraint_Error\| \|C/C++\|\|modulo power of two\|\|undefined behavior\| \|C#\|\|modulo power of 2 in unchecked context; \|Python 2\|\|NA\|\|convert to long\| \|Scheme\|\|NA\|\|convert to bigNum\| \|Smalltalk\|\|NA\|\|convert to LargeInteger\| \|Swift\|\|Causes error unless using special overflow operators.\| There are several methods of handling overflow: - Avoidance: by carefully ordering operations, checking operands in advance and selecting the correct data type, it is possible to ensure that the result will never be larger than can be stored. - Handling: If it is anticipated that overflow may occur and when it happens detected and other processing done. Example: it is possible to add two numbers each two bytes wide using just a byte addition in steps: first add the low bytes then add the high bytes, but if it is necessary to carry out of the low bytes this is arithmetic overflow of the byte addition and it necessary to detect and increment the sum of the high bytes. CPUs generally have a way of detecting this to support addition of numbers larger than their register size, typically using a status bit. - Propagation: if a value is too large to be stored it can be assigned a special value indicating that overflow has occurred and then have all successive operation return this flag value. This is useful so that the problem can be checked for once at the end of a long calculation rather than after each step. This is often supported in Floating Point Hardware called FPUs. Programming languages implement various mitigation methods against an accidental overflow: Ada, Seed7 (and certain variants of functional languages), trigger an exception condition on overflow, while Python (since 2.4) seamlessly converts internal representation of the number to match its growth, eventually representing it as long—whose ability is only limited by the available memory. Run-time overflow detection implementation AddressSanitizer is also available for C compilers. In languages with native support for Arbitrary-precision arithmetic and type safety (such as Python or Common Lisp), numbers are promoted to a larger size automatically when overflows occur, or exceptions thrown (conditions signaled) when a range constraint exists. Using such languages may thus be helpful to mitigate this issue. However, in some such languages, situations are still possible where an integer overflow can occur. An example is explicit optimization of a code path which is considered a bottleneck by the profiler. In the case of Common Lisp, this is possible by using an explicit declaration to type-annotate a variable to a machine-size word (fixnum) and lower the type safety level to zero for a particular code block. In Java 8, there are overloaded methods, for example like Math#addExact(), which will throw ArithmeticException in case of overflow. Computer emergency response team (CERT) developed the As-if Infinitely Ranged (AIR) integer model, a largely automated mechanism to eliminate integer overflow and truncation in C/C++ using run-time error handling. In computer graphics or signal processing, it is typical to work on data that ranges from 0 to 1 or from −1 to 1. An example of this is a grayscale image where 0 represents black, 1 represents white, and values in-between represent varying shades of gray. One operation that one may want to support is brightening the image by multiplying every pixel by a constant. Saturated arithmetic allows one to just blindly multiply every pixel by that constant without worrying about overflow by just sticking to a reasonable outcome that all these pixels larger than 1 (i.e. "brighter than white") just become white and all values "darker than black" just become black. Unanticipated arithmetic overflow is a fairly common cause of program errors. Such overflow bugs may be hard to discover and diagnose because they may manifest themselves only for very large input data sets, which are less likely to be used in validation tests. Taking the arithmetic mean of two numbers by adding them and dividing by two, as done in many search algorithms, causes error if the sum (although not the resulting mean) is too large to be represented, and hence overflows. An unhandled arithmetic overflow in the engine steering software was the primary cause of the crash of the 1996 maiden flight of the Ariane 5 rocket. The software had been considered bug-free since it had been used in many previous flights, but those used smaller rockets which generated lower acceleration than Ariane 5. On 30 April 2015, the Federal Aviation Authority announced it will order Boeing 787 operators to reset its electrical system periodically, to avoid an integer overflow which could lead to loss of electrical power and ram air turbine deployment, and Boeing is going to deploy a software update in the fourth quarter. The European Aviation Safety Agency followed on 4 May 2015. The error happens after 2³¹ centiseconds (248.55134814815 days), indicating a 32-bit signed integer. Overflow bugs are evident in computer games. In the arcade game Donkey Kong, it is impossible to advance past level 22 due to an integer overflow in its time/bonus. The game takes the level number a user is on, multiplies it by 10 and adds 40. When they reach level 22, the time/bonus number is 260, which is too large for its 8-bit 256 value register, so it resets itself to 0 and gives the remaining 4 as the time/bonus - too short to finish the level. In Donkey Kong Jr. Math, when trying to calculate a number over 10000, it shows only the first 4 digits. Overflaw is the cause of the famous Split Screen in Pac-Man and the Nuclear Gandhi in Civilization series. - Buffer overflow - Heap overflow - Pointer swizzling - Software testing - Stack buffer overflow - Static program analysis - Unix signal - Seed7 manual, section 15.2.3 OVERFLOW_ERROR. - The Swift Programming Language. Swift 2.1 Edition. October 21, 2015. - Python documentation, section 5.1 Arithmetic conversions. - Reddy, Abhishek (2008-08-22). "Features of Common Lisp". - Pierce, Benjamin C. (2002). Types and Programming Languages. MIT Press. ISBN 0-262-16209-1. - Wright, Andrew K.; Matthias Felleisen (1994). "A Syntactic Approach to Type Soundness". Information and Computation. 115 (1): 38–94. doi:10.1006/inco.1994.1093. - Macrakis, Stavros (April 1982). "Safety and power" (requires subscription). ACM SIGSOFT Software Engineering Notes. 7 (2): 25–26. doi:10.1145/1005937.1005941. - As-if Infinitely Ranged Integer Model - Google Research blog: Nearly All Binary Searches and Mergesorts are Broken, Joshua Bloch, 2 June 2006 - Gleick, James (1 December 1996). "A Bug and A Crash". New York Times Magazine. Retrieved 9 December 2013. - "F.A.A. Orders Fix for Possible Power Loss in Boeing 787". New York Times. 30 April 2015. - "US-2015-09-07 : Electrical Power - Deactivation". Airworthiness Directives. European Aviation Safety Agency. 4 May 2015. - Pittman, Jamey. "The Pac-Man Dossier".
Principles of Finance/History Introduction to Finance Finance is a field of study of the relationship of three things; time, risk and money. The Time Value of Money is one of three fundamental ideas that shape finance. The Time Value of Money explains why, "A dollar today is worth more than a dollar tomorrow". This is primarily due to the market for loanable funds and inflation. If someone has a dollar today then they also have the opportunity to loan/invest that dollar at some interest rate. Therefore, a dollar today in time t, would be worth $1.00 plus some interest rate, i. That is more than a dollar by itself in the future. An example for inflation would be, let's say you have $1 and you can buy 10 candies today. For the same 10 candies tomorrow you have to pay $1.20. So, due to inflation, for the same 10 candies, today you pay less than that you would pay tomorrow Inflation refers to the decrease in the purchasing power. Deflation refers to the increase in the purchasing power. In layman terms, inflation causes not the value of money to decrease but the amount of consumables/items that you can now purchase to decline in quantity. Look at the example above. $1 is still $1 but after inflation the individual can probably buy only 8 candies for the same $1 amount. There are two values of money. One is the Present Value of Money and the other is the Future Value of Money. Second is the concept of "opportunity cost"; i.e. if a person deploys his money on one item or investment then he has given up the opportunity to do something else with it. Third is the concept of risk. Let us say, I have earmarked $10,000 towards investments and I decide that I will invest in Microsoft. I put all my money in Microsoft(MSFT) all $10,000 of it. As on 5th Oct., 2006 each share of Microsoft trades at $27.94. So, I would be able to purchase 357.91 shares of MSFT. My returns are completely dependent on how MSFT stock performs in the market and this means that if a Microsoft product(i.e. Vista), fails in the marketplace, MSFT stock goes tumbling down and reduces my investment in MSFT. Thus, Risk can be defined as the probability of my investment eroding its own self. In equity, the risk factor is higher than in debt financing and hence as an investor I look for equity to give me higher returns as I have taken higher risk. If I buy US Treasury notes for that value, the investment is almost risk free as the US Govt. stands to guarantee it and hence the return (typically, expressed as the rate of interest) is low. The difference between the returns from equity and from debt(US Treasury, etc) is the Risk Premium. Risk Premium is the reward given to an investor to take more risk. Return on Investments The concept of a return on investment is designed to balance all three perils. In finance there are actually two returns: the return of investment and the return on investment. If the investment loses money, you may be able to recover a portion of what you risked. A return on investment (hereafter "ROI", or "return") is the return that finance is primarily concerned with although the other cannot be overlooked. It implies first and foremost a 100% return of investment. In order to do so, the return must therefore - replace the buying power lost to inflation, - make up for and exceed the losses from other financing activities, and - make the investment more attractive to someone than any other option, including spending the money. Item 3 can be as objective as selecting the best ROI among many or it can be as subjective and personal as whether or not to give up the satisfaction of having dessert every night for a month. Regardless if the investor does not perceive sufficient potential for gain, the money will never change hands. Debt Finance and Equity Finance - The Two Pillars of Modern Finance Financing activity for most ventures are either debt financing or equity financing. Once an investor has decided to engage in financing any business venture or project, he has three concerns to address: risk, protective claim, and return. These three concerns are essentially hierarchical. The latter two depend on and counterbalance with the first. Higher risk is only attractive when associated with higher returns. The protective claim then acts like a fulcrum to fine tune the balance between risk and return. Debt financing. In this mode, money is borrowed, and usually the borrower (debtor) gives the lender (creditor) a promisory note. This, usually, obligated the debtor to pay back a certain defined amount at a particular and defined time in the future. Forms of debt financing can include credit cards, mortgages, signature loans, bonds, IOUs, and HELOCs (Home Equity Lines of Credit) as examples. Treasury debt, savings bonds, corporate bonds are also forms of debt financing. With debt financing, the creditor's return is fixed and understandable. It is, quite simply, the agreed upon interest rate for the debt. This rate can vary from a single digit rate to 25% or perhaps even 30% depending on the debtor. Risk is determined by a handful of factors the most significant usually being one's credit history. The protective claim offered to creditors in debt financing is a claim on the debtor's assets. Should the debtor fail to repay, the creditor may forcibly take possession of other debtor property and sell it, using the money to offset the loss of the loan. The claim of creditors takes priority over the claim of those who participate in equity financing. Equity financing is generally considered less certain than debt financing. Equity financing is also typically where non-cash assets such as equipment, skills, and land are invested alongside regular cash. This is the category in which we also find venture capital, shares of stock, angel investors, and more. The terms that are used to describe the equity financing relationship are more varied and, as such, will be simply dubbed equity investors. Later on, more proper names will be provided which, for the time being, are immaterial. The return of equity financing is the claim on a business's profits; not just today's profits, but in modern companies that issue stock, all future potential profits as well. For this reason, i.e. because most personal finance does not involve the debtor making a profit, almost all of personal finance is debt financing. The exceptions will be noted shortly. While it's true that in equity financing, the equity investor still has some claim on the business' assets, the creditor's prior claim renders this point moot from a practical standpoint. What protection, then, is offered for equity financing? The claim to management rights. As an equity investor, with a few notable exceptions, you are granted the right to do everything in your personal power along with the other equity investors to make sure that the business goes in a profitable direction. Ratio analysis is one core theme of business finance. To understand the concept of ratio analysis we must know that there are five basic types of ratio analysis. - Liquidity Ratio - Activity Ratio - Debt Ratio - Profitability Ratio - Markets Ratio. Now we discuss one by one the impacts of these ratios on the business and their implementation in the business environment. In liquidity ratio analysis, there are two types of ratios: - Current Ratio - Quick (Acid-Test) Ratio. Current Ratio shows how many $1 of current assets are available for paying $1 of current liabilities of the company, firm or organization. As some financial analysts suggests that the more current assets that are available to the company, the better. But in some cases, a high current ratio may indicate too much inventory or too much in prepaid or other current assets. It could also indicate idle cash, and as a result, a poor investment strategy. A current ratio that is high is not as bad as one that is low, however a high current ratio is an indication that financial policies either do not exist or are not being implemented. A low or declining current ratio is not always bad. A declining ratio is an indicator of rising current liabilities and declining current assets. But a current ratio of 1 or less is an indication of insolvency. Further declines in the ratio could trigger collection actions on the part of creditors and send the company into bankruptcy. Quick (Acid-Test) Ratio The quick ratio is calculated by dividing cash by current liabilities. This is the true test of a companies ability to pay its debts. A high current ratio and a low quick ratio could indicate too much is invested in inventory or other current assets, assets that are not very liquid and therefore could not be depended on to pay current liabilities. An increase in inventory could be a signal that sales have fallen, production has slowed and management should take action to prevent any damage to the financial condition of the company. Both ratios should be analyzed together to get the correct picture of the companies financial health.
We are searching data for your request: Upon completion, a link will appear to access the found materials. William Penn became a member of the Society of Friends, otherwise known as the Quakers. They believed in a simple life style. They believed that all men were equal. Quakers refused to bow to the King or to fight in wars. They also refused to pay taxes to church. King Charles II tried to stop Penn from preaching in favor of Quaker beliefs by briefly imprisoning him, but that did not stop Penn from continuing to preach. Biography of Penn Dutch and Swedes Arrive The Indians of Pennsylvania — [Four] hundred years ago the region now known as Pennsylvania had never felt the tread of a [European person's] foot. White settlers had come to other parts of the country, but here dwelt only the [Native Americans], those natives of the land whom we call Indians. Chief among these were those known as Delawares, from the river on which they dwelt, but who called themselves the Lenni- Lenapes. The tribe of the Delawares was divided into three sections, or sub-tribes, the Minsi, or Minisink, the Unami, and the Unalachtigo, which had respectively for totems the wolf, the turtle, and the turkey. The Unami, or Turtle, section dwelt on the site of Philadelphia. Other tribes, separate from the Delawares, were the Susquehannocks, the Nanticokes, and the Eries, who lived farther west. The Peaceful Delawares — The [European] settlers of Pennsylvania had most to do with the Delawares, who, by good fortune, were a peaceful people. They had been conquered by the warlike Iroquois of New York and forced by them to keep peace with all the tribes. Instead of making war they were to till the soil as women did, and to them was given the care of "the great belt of peace." At a later date another tribe, called the Shawnees, came to Pennsylvania, a few of them at first, but eventually there were many of them in the province. Such were the native tribes found by William Penn and his Quaker friends when they crossed the ocean to America. Visitors before the Quakers — The Quakers were not the first [Europeans] to reach Pennsvlvania. Others were there before them. When we speak of how this province was settled we are apt to think first of William Penn, but long before he came many settlers had reached this locality. The history of these early settlers must be told before we speak of Penn. There were Swedes, Dutch, and English, about each of whom there is something to tell. The first man to sail up the Delaware was a Dutch captain named Hendrickson, who in 1616 went up this fine river as far as the mouth of the Schuylkill. He was much pleased with what he saw there, for he had found a beautiful land, with a great forest full of deer, turkeys and partridges, and with vines clambering up the trees. There was also a Captain Mey, from whom Cape May got its name, who in 1623 sailed up the river and built a fort at a point four miles below the site of Philadelphia. This he named Fort Nassau. In 1630 a small party of Dutch settled near the lower end of Delaware. But a foolish quarrel soon put an end to their settlement. They had painted the arms of Holland on a piece of tin and hung it up on a tree. An Indian took it down to make a tobacco pipe, and for this he was killed, either by the Dutch or by the members of his tribe in consequence of the angry protests of the settlers, to whom the act of the ignorant native, who knew nothing about the arms of Holland, seemed an insult to their country. The death of the Indian was quickly avenged by his friends, who attacked the settlement and killed every person in it. Thus ended in crime and blood the first settlement on the Delaware [River]. The Coming of the Swedes — It was not long before new settlers came. In 1637 two small vessels set sail from Sweden [arriving in 1638], loaded with Swedes and Finns, who sought a new home on the banks of the South River—as the Dutch called the Delaware. They were led by Peter Minuit, a Dutchman, who knew the country well, for he had been governor of the Dutch settlement of New Amsterdam. He bought from the Indians all the land on the west shore of the Delaware as far up as the mouth of the Schuylkill, built a fort where Wilmington now stands, and named it Fort Christina, in honor of the Queen of Sweden. A new governor, named Hollender, came in 1641, and bought from the Indians a large tract of land along the river, and in 1643 there came a third governor, named Johan Printz, who built himself a fine mansion and a strong fort on Tinicum Island, a few miles below Philadelphia, and lived there in much style. The Swedes called their colony New Sweden and claimed all the land on the west side of the Delaware from Cape Henlopen to Trenton Falls. They also claimed the east side from Cape May to Mantua Creek, nearly opposite Chester. They traded for furs with the Indians, planted wheat, rye, and tobacco, and built forts for defense. The End of New Sweden — Bv 1650 the Swedes had a thriving settlement. Much land was cleared and planted, they had plenty of fruit, grain, and cattle, and built a mill on Cobb's Creek, which was kept busy grinding their grain. But the Dutch of New Amsterdam had been first on the ground, had built forts and bought land from the Indians, and though they had not settled the country they did not like to see the wav the Swedish Colony was growing. So they collected a little fleet with an army of about six hundred men and in August, 1655, set sail for the [Delaware] River. This was not a very large army, but the Swedes, not being strong enough to fight, gave up to the Dutch without firing a shot or striking a blow. They were left on their farms under the rule of Holland and the colony of New Sweden came to an end. Relics of New Sweden — The settlements of the Swedes lay along the west side of the river from New Castle, in Delaware, to the site of Philadelphia. They had built a church on Tinicum Island in 1646, and a church was built about 1669 at Wicaco in what is now southern Philadelphia. This was rebuilt in later years, and still stands, known as the Gloria Dei, or Old Swedes Church. They had a small town at Upland—now Chester—and here their first courts were held, the first jury sat, and the first highway was built. The English Claim — New changes were soon to come, for the English also claimed this region. In 1664 an English fleet appeared before New Amsterdam, the Dutch settlement on Manhattan Island, took it without firing a gun, and named it New York. Then they sent two ships to the Delaware and took the settlement there also, but not until some Dutch soldiers had been killed and wounded. This was the first bloodshed in all the quarrels of the [Europeans] in that region. The Swedes were quite willing to come under English rule, and so were the Dutch, for they were well treated by their new masters, their farms left in their hands, and all their officers left in their posts. There were not many of them, probably only a few hundred in all, and they were widely scattered along the river. New Castle was the centre of government and Upland the place of next importance. Philadelphia was still only a region of farms. The Quaker colony In March 1681 Charles II of England signed a charter giving any unoccupied regions to William Penn in payment of a debt owed by the king to Penn’s father, Adm. Sir William Penn. The charter, which was officially proclaimed on April 2, 1681, named the territory for Admiral Penn and included also the term sylvania (“woodlands”), at the son’s request. William Penn intended that the colony provide a home for his fellow Quakers (members of the Society of Friends). While still in England, he drew up the first of his “frames of government” and sent his cousin, William Markham, to establish a claim to the land and also to establish the boundaries of what became the city of Philadelphia. Penn arrived in 1682 and called a General Assembly to discuss the first Frame of Government and to adopt the Great Law, which guaranteed freedom of conscience in the colony. Under Penn’s influence, fair treatment was accorded the Native Americans, who responded with friendship in return. When Penn returned to England in 1684, the new Quaker province had a firmly established government based on the people’s will and religious tolerance. Early History of Native Americans in Pennsylvania The names of the Pennsylvania tribes included the Lenapi Delaware, Erie, Honniasont, Iroquois, Saponi, Shawnee, Susquehanna, Tuscarora, Tutelo and Wenrohronon. Native Americans lived in the area that became Pennsylvania hundreds of years before European settlers entered the region. The two primary groups were the Algonkian and Iroquois. Algonkian tribes included the Delaware, Nanticoke, and Shawnee. The Susquehannocks were an Iroquoian tribe that lived along the Susquehanna River. These early inhabitants traveled by canoe or on foot. They lived in houses made of bark and wore clothing from the skins of animals. Arts such as pottery making and weaving were also practiced. Although some farming was done, most food was acquired through hunting and gathering. When first discovered by Europeans, Pennsylvania, like the rest of the continent, was inhabited by groups of American Indians, people of Mongoloid ancestry unaware of European culture. The life of the Indians reflected Stone Age backgrounds, especially in material arts and crafts. Tools, weapons and household equipment were made from stone, wood, and bark. Transportation was on foot or by canoe. Houses were made of bark, clothing from the skins of animals. The rudiments of a more complex civilization were at hand in the arts of weaving, pottery, and agriculture, although hunting and food gathering prevailed. Some Indians formed confederacies such as the League of the Five Nations, which was made up of certain New York-Pennsylvania groups of Iroquoian speech. The other large linguistic group in Pennsylvania was the Algonkian, represented by the Delawares, Shawnees, and other tribes. The Six Nations of the Iroquois Confederacy and the Seneca Nation occupied the lands now known as Erie. For a history of the Native Americans who occupied this land before Europeans, see Erie Indians. The French built Fort Presque Isle near present-day Erie in 1753, as part of their effort to garrison New France against the encroaching English. The French word "Presque-isle" means peninsula (literally "almost an island") and refers to that piece of land that juts into Lake Erie that is now called Presque Isle State Park. When the fort was abandoned by the French in 1760, it was their last post west of Niagara. The British occupied the fort at Presque Isle that same year, three years before the end of the Seven Years' War in 1763. Present day Erie would have been situated in a disputed triangle of land that was claimed by the states of New York, Pennsylvania, Connecticut (as part of its Western Reserve), and Massachusetts. It officially became part of Pennsylvania on 3 March 1792, after Connecticut, Massachusetts, and New York released their claims to the federal government, which in turn sold the land to Pennsylvania for $151,600 in Continental certificates. The Six Nations of the Iroquois Confederacy released the land to Pennsylvania in January 1789 for payments of $2,000 from Pennsylvania and $1,200 from the federal government. The Seneca Nation separately settled land claims against Pennsylvania in February 1791 for the sum of $800. The General Assembly of Pennsylvania commissioned the surveying of land near Presque Isle through an act passed on 18 April 1795. Andrew Ellicott, who famously completed Pierre Charles L'Enfant's survey of Washington, D.C. and helped resolve the boundary between Pennsylvania and New York, arrived to begin the survey in June 1795. Initial settlement of the area began that year. In 1795, Colonel Seth Reed and his family, natives of Uxbridge, Massachusetts, moved here from Geneva, New York, to become the first European settlers of Erie. Reed erected a log cabin at the mouth of Mill Creek, becoming the first permanent building in Erie. Reed's other sons, Rufus S. Reed and George W. Reed, came to Erie later in the year. Erie was established as a borough by act of the General Assembly on 29 March 1805. This act created a Borough and Town Council headed by a burgess. This form of government stood until the City of Erie was incorporated on 14 April 1851, when a mayoralty and Select Council were established. During the War of 1812, President James Madison ordered the construction of a naval fleet at Erie in order to regain control of Lake Erie. Noted shipbuilders Daniel Dobbins of Erie and Noah Brown of New York led construction of four schooner-rigged gunboats and two brigs. Commodore Oliver Hazard Perry arrived from Rhode Island to command the squadron. His fleet successfully fought the British in the historic Battle of Lake Erie, which was the decisive victory that solidified United States control of the Great Lakes. Erie was an important railroad hub in the mid-19th century. However, the railroad north to Buffalo, New York used 6' track up to the New York border, while the railroad west to Cleveland, Ohio and that from the New York border to Buffalo from was on the narrower 4' 10" gauge. What this meant was that there was no line through Erie every passenger had to change trains, and every piece of cargo had to be moved by railroad stevedores and wagons between trains. While the trains were timed to connect, delays could cause passengers to miss their connection they then needed a meal and a bed in Erie. The delays inconvenienced both passengers and cargo, adding to the time and therefore the expense of travel by rail between Buffalo to Cleveland. However, they provided much needed jobs in Erie. Travelers were patrons of Erie hotels, restaurants, and stores. Those shipping goods needed manpower, and some of this came from Erie itself there were many self-employed "men with a horse" and a wagon moving goods. The two railroads themselves provided jobs. It was obvious to everyone except those from Erie that this was a ridiculous situation. The 6' section needed to be changed 4'10", the national standard, so trains could go through Erie, and passengers and goods would not have to change trains twice between Buffalo and Cleveland. However, this would have a substantial negative effect on employment and the economy in Erie, which benefited from the unavoidable train change. The citizens of Erie, led by the mayor, set fire to bridges, ripped up track, and in general did everything imaginable to prevent the change. Erie's congressional representative Milton W. Shreve supported the Volstead Act and the Eighteenth Amendment. Miles Nason, another Erie Prohibitionist, headed the Dry Block in the Pennsylvania State Senate. But Erie was primarily a "wet" city. Being a border town, Erie was an important transportation hub in the rum-running of illicit liquor across the lake from Canada during Prohibition in the United States. John G. Carney, in his "Highlights of Erie Politics", says that many "laid in a large supply of liquor before the law became effective. Cellars, book cases, and closets were packed. " Speakeasies opened across the city, the more popular being the Pickwick Club, the Killarney Yacht Club, Laura's, and 1008. Carney noted that ". about the only dry thing in Erie was the inside of a light bulb." Illicit liquor sales brought racketeering, violence, and houses of prostitution. Intervention by the state police was not welcomed by Mayor Miles B. Kitts, who went to Harrisburg and testified before well-publicized hearings conducted by Pennsylvania Governor William C. Sproul. But the actions of local and state law enforcement and the governor's hearings offered only a brief respite from all the excitement. As Carney concluded, ". and Erie 'roared' merrily on throughout the rest of the 'Roaring Twenties.' " Shreve fell from favor with the Republicans, who promoted attorney Robert Firman as their candidate in the April 1920 primaries. Shreve narrowly escaped removal from the United States Congress. State Senator Nason was also challenged by the Republicans in the primaries, but was defeated in the 2 November 1920 elections. The Great Depression deflated Erie's enthusiasm for lawlessness and prompted a solid political movement towards repeal of Prohibition. Democratic Party chairman for Erie County and future mayor James P. Rossiter was able to promise strong local voter support for Democratic-Liberal candidate for state governor John Hemphill when he visited Erie with a strong agenda for repeal in October 1930. In 2007, the Erie Downtown Improvement District (DID) contracted a Philadelphia-based company (Kise, Straw, & Kolodner) to set up a "master plan" for the city of Erie's downtown. The DID plan includes building several mid-rise and high-rise structures which will be used primarily for housing and retail expansion in the city center. Fourth River Development and Radnor Property Group were selected as the developers. In January 2007, GAF, an asphalt shingle manufacturer announced plans to relocate to Eastern Pennsylvania, thus making available several extremely valuable acres next to the Convention Center and hotel under construction. A local newspaper poll showed that the majority of local citizens desire a park-like setting, followed by retail development in the area. Our editors will review what you’ve submitted and determine whether to revise the article. Pennsylvania German, also called (misleadingly) Pennsylvania Dutch, 17th- and 18th-century German-speaking settlers in Pennsylvania and their descendants. Emigrating from southern Germany (Palatinate, Bavaria, Saxony, etc.) and Switzerland, they settled primarily in the southeastern section of Pennsylvania, where they practiced any of several slightly different forms of Anabaptist faith, mostly Amish and Mennonite. Their descendants, some of whom participate only reluctantly in modern life, live mainly in Northampton, Berks, Lancaster, Lehigh, Montgomery, Bucks, York, and other counties of Pennsylvania, as well as in Ohio, Indiana, Iowa, Kansas, Oklahoma, Virginia, West Virginia, and Florida. Some groups—especially those who remain apart—still speak (in addition to English) a German dialect known as Pennsylvania Dutch or Pennsylvania German, a blending of High German (in reference to the altitude of their natal region), various German dialects, and English. The word Dutch (from German Deutsch, meaning “German”), which once encompassed all non-English speakers of Germanic languages, is in the 21st century a misnomer, as Dutch has come to be associated strictly with people from the Netherlands. Many Pennsylvania Germans are thoroughly assimilated, though they may retain elements of their traditional culture such as special cookery (e.g., shoofly pie, an extremely sweet pie made with molasses and brown sugar) and a decorative tradition known as fraktur (which blends calligraphic and pictorial elements). Some groups, such as the Old Order Amish, wear plain, modest clothing and head coverings and drive horse-drawn buggies. Men wear beards (but not mustaches) after they marry. They live according to relatively strict religious principles. The Pennsylvania Germans, many of whom had been persecuted in their native land, were attracted to Pennsylvania by the liberal and tolerant principles of William Penn’s government. Their immigration began with the Mennonite Francis Daniel Pastorius, who in 1683 led a group of German Quakers to Philadelphia, where they founded Germantown, the pioneer German settlement. The early German settlers were for the most part Mennonites, Amish, Dunkers (or German Baptists), Schwenckfelders, and Moravians (see Moravian church). After 1727 the immigrants were mostly members of the larger Lutheran and Reformed churches. Their farming skills made their region of settlement a rich agricultural area. By the time of the American Revolution they numbered about 100,000, more than a third of Pennsylvania’s population. Philadelphia, a city in Pennsylvania whose name means City of Brotherly Love, was originally settled by Native American tribes, particularly the Lenape hunter gatherers, around 8000 B.C. By the early 1600s, Dutch, English and Swedish merchants had established trading posts in the Delaware Valley area, and in 1681, Charles II of England granted a charter to William Penn for what would become the Pennsylvania colony. Penn arrived in the new city of Philadelphia in 1682. A Quaker pacifist, Penn signed a peace treaty with Lenape chief Tamanend, establishing a tradition of tolerance and human rights. But in 1684, the ship Isabella landed in Philadelphia carrying hundreds of enslaved Africans. Tensions over slavery, especially among local Quakers, resulted in the 1688 Germantown Petition Against Slavery, the first organized protest against slavery in the New World. Penn’s colony thrived, and soon Philadelphia was the biggest shipbuilding center in the colonies. Among those attracted to the city was Benjamin Franklin, who in 1729, became the publisher of The Pennsylvania Gazette. The Pennsylvania State House—later known as Independence Hall—held its first Assembly meeting there in 1735. State representatives ordered a large bell for the building in 1751 with a Biblical inscription: “Proclaim LIBERTY throughout all the Land unto all the inhabitants thereof.” British Parliament passed a series of tax acts on the colonies in the 1760s, including the Stamp Act and the Townshend Acts, sparking colonial outrage. In response, the Continental Congress convened in Philadelphia in 1774. After Philadelphia resident Thomas Paine&aposs pamphlet Common Sense met with widespread acclaim, the stage was set to formally declare independence, which the Founding Fathers did on July 4, 1776. Philadelphians were the first to hear the Declaration of Independence read aloud in the State House yard. In 1790, after the Revolutionary War (during which the city witnessed the Battle of Germantown), Philadelphia served as capital of the United States. By that time, it was the new nation’s biggest city, with 44,096 residents. The First Bank of the United States and the first U.S. Mint were founded in Philadelphia, and the U.S. Constitution was written there in 1787. With the city’s history of civil rights—the Pennsylvania Abolition Society met there in 1775— Philadelphia was an ideal spot for William Lloyd Garrison to establish the American Anti-Slavery Society, which grew to nearly 250,000 members by 1838. Local abolitionists adopted the old State House bell as a symbol, renaming it the “Liberty Bell.” Philadelphia rallied to the Union cause during the Civil War, and local industries profited by supplying weapons, uniforms and warships. In 1876, suffragette Susan B. Anthony delivered the Declaration of the Rights of Women outside Independence Hall. The city grew in size and prestige during the Gilded Age, as wealthy suburbs sprouted along the Main Line of the Pennsylvania Railroad. During the 1870s, the first U.S. zoo and the Centennial Exhibition fair opened in Philadelphia. The city’s shipbuilding industries supplied the Allies in World War I, but Philadelphia was also a center of the Spanish flu pandemic of 1918-1919—over 500,000 citizens contracted the deadly disease. After World War II, new highways allowed workers to easily reach bedroom communities outside the city. With suburbanization and industrial decline, Philadelphia lost population and jobs, and soon many of the city’s famed shipyards were shuttered. Poverty and racial tensions soon followed, and in 1985 a police confrontation with the radical group MOVE ended with the bombing of a predominantly black neighborhood people in the MOVE compound were killed. New developments, such as the Philadelphia Navy Yard and Center City, have helped to revitalize the area, which is now home to more than 1.5 million residents. The city rejoiced when the Eagles won the 2018 Super Bowl. For visitors, a perennially popular destination is the statue of Rocky Balboa, depicting the fictional boxer, arms outstretched, at the top of the steps to the Philadelphia Art Museum. Rocky, played by Sylvester Stallone, famously runs up the 72 steps to train for a fight in the 1976 movie, "Rocky" (and in sequels). Now the stairs to the museum are simply known as the "Rocky Steps." In 1681, King Charles II gave Penn a large piece of his newly acquired American land holdings to repay a debt the king owed to Admiral Sir William Penn, Penn's father. This land included present-day Pennsylvania and Delaware, though the claim as written would create a bloody conflict with Maryland (dubbed Cresap's War) over the land grant already owned by Lord Baltimore. Penn put together a colonial expedition and fleet, which set out for America in the middle of the following summer. Penn, sailing in the vanguard, first set foot on American soil at the colony at New Castle, Delaware. An orderly change of government ensued, as was normal in an age used to the privileges and prerogatives of aristocracy and which antedated nationalism: the colonists pledged allegiance to Penn as their new Proprietor. The first Pennsylvania General Assembly was soon held in the colony. Afterwards, Penn journeyed up the river and founded Philadelphia with a core group of accompanying Quakers and others seeking religious freedom on lands he purchased from the local chieftains of the Lenape or Delaware nation. This began a long period of peaceful co-operation between the colony and the Delaware, in contrast to the frictions between the tribe and the Swedish and Dutch colonists. [ page needed ] However, the new colonists would not enjoy such easy relations with the rival and territorial Conestoga peoples to the west for a number of decades as the English Quaker and German Anabaptist, Lutheran and Moravian settlers attracted to the religiously tolerant colony worked their way northwest up the Schuylkill and due west south of the hill country into the breadbasket lands along the lower Susquehanna River. Lord Baltimore and the Province of Maryland had circa 1652–53 finished waging a decade long declared war against the Susquehannocks and the Dutch, who'd been trading them furs for tools and firearms for some time. Both groups had uneasy relations with the Delaware (Lenape) and the Iroquois. Furthermore, Penn's Quaker government was not viewed favorably by the Dutch, Swedish, and English settlers in what is now Delaware. They had no "historical" allegiance to Pennsylvania, so they almost immediately began petitioning for their own Assembly. In 1704, they achieved their goal when the three southernmost counties of Pennsylvania were permitted to split off and become the new semi-autonomous colony of Lower Delaware. New Castle, the most prominent, prosperous and influential settlement in the new colony, became the capital. During its brief period of ascendancy as an empire following the victory by Gustav the Great in the Battle of Breitenfeld Swedish settlers arrived in the area in the early 17th century to found a nearby colony, New Sweden in what is today southern New Jersey. With the arrival of more numerous English colonists and development of the port on the Delaware, Philadelphia quickly grew into an important colonial city. During the American Revolution, it was the site of the First and Second Continental Congresses. After the Revolution, the city was chosen to be the temporary capital of the United States from 1790 to 1800. At the beginning of the 19th century, the federal and state governments left Philadelphia, but the city continued for some years to be the country's cultural and financial center. Its large free black community aided fugitive slaves and founded the first independent black denomination in the nation, the African Methodist Episcopal Church. Philadelphia became one of the first U.S. industrial centers with a variety of industries, the largest being textiles. It had many economic and family ties to the South, with southern planters maintaining second homes in the city and having business connections with banks, sending their daughters to French finishing schools run by refugees from Saint-Domingue (Haiti), selling their cotton to textile manufacturers, which in turn sold some products to the South, for instance, clothing for slaves. At the beginning of the American Civil War, there were many southern sympathizers, although most city residents became firmly Union as the war went on. After the American Civil War, city government was controlled by the Republican Party it established a political machine that gained power through patronage. By the beginning of the 20th century, Philadelphia was described as "corrupt and contented." Various reform efforts slowly changed city government in 1950, a new city charter strengthened the position of mayor and weakened the Philadelphia City Council. Beginning during the Great Depression, voters changed from traditional support for the Republican Party to increasing support for the Democratic Party of President Franklin D. Roosevelt, which has now been predominant in local politics for many decades. The population grew dramatically at the end of the 19th and beginning of the 20th centuries, through immigration from Ireland, Southern Europe, Eastern Europe, and Asia, as well as the Great Migration of blacks from the rural South and Puerto Ricans from the Caribbean, all attracted to the city's expanding industrial jobs. The Pennsylvania Railroad was expanding and hired 10,000 workers from the South. Manufacturing plants and the US Navy Yard employed tens of thousands of industrial workers along the rivers, and the city was also a center of finance and publishing, with major universities. By the 1950s, much Philadelphia housing was aged and substandard. In the post-World War II era of suburbanization and construction of area highways, many middle-class families met their demand for newer housing by leaving the city for the suburbs. Population decline accompanied the industrial restructuring and the loss of tens of thousands of jobs in the mid 20th century. With increasing poverty and social dislocation in the city, gang and mafia warfare plagued the city in from the mid-20th century to the early 21st century. By the end of the 20th century and beginning of the 21st, revitalization and gentrification of historic neighborhoods attracted an increase in middle-class population as people began to return to the city. New immigrants from Southeast Asia, and Central and South America have contributed their energy to the city. Promotions and incentives in the 1990s and the early 21st century have improved the city's image and created a condominium boom in Center City and the surrounding areas. Before Philadelphia was colonized by Europeans, the area was inhabited by the Lenape (Delaware) Indians. The village of Nitapèkunk, "place that is easy to get to," was located in today's Fairmount Park area. The villages of Pèmikpeka, "where the water flows," and Shackamaxon were located on the Delaware River. The Delaware River Valley was called the Zuyd, meaning "South" River, or Lënapei Sipu. The first exploration of the area by Europeans was in 1609, when a Dutch expedition led by Henry Hudson entered the Delaware River Valley in search of the Northwest Passage. The Valley, including the future location of Philadelphia, became part of the New Netherland claim of the Dutch. Dutch explorer Cornelius Jacobsen May charted the shoals Delaware Bay in the 1620s and a fort was built on the west side of the bay at Swanendael. In 1637, Swedish, Dutch and German stockholders formed the New Sweden Company to trade for furs and tobacco in North America. Under the command of Peter Minuit, the company's first expedition sailed from Sweden late in 1637 in two ships, Kalmar Nyckel and Fogel Gri. Minuit had been the governor of the New Netherland from 1626 to 1631. Resenting his dismissal by the Dutch West India Company, he brought to the new project the knowledge that the Dutch colony had temporarily abandoned its efforts in the Delaware Valley to focus on the Hudson River valley to the north. (The Hudson was known to the Dutch as the Noort, or "North" river relative to "South" of the Delaware.) Minuit and his partners also knew that the Dutch view of colonies necessitated occupation to secure legal claim. The ships reached Delaware Bay in March 1638, and the settlers began to build a fort at the site of present-day Wilmington, Delaware. They named it Fort Christina, in honor of the twelve-year-old Queen Christina of Sweden. It was the first permanent European settlement in the Delaware Valley. Part of this colony eventually included land on the west side of the Delaware River from just below the Schuylkill River. [ citation needed ] Johan Björnsson Printz was appointed to be the first royal governor of New Sweden, arriving in the colony on February 15, 1643. Under his ten-year rule, the administrative center of New Sweden was moved north to Tinicum Island (to the immediate SW of today's Philadelphia), where he built Fort New Gothenburg and his own manor house which he called the Printzhof. The first English settlement occurred about 1642, when 50 Puritan families from the New Haven Colony in Connecticut, led by George Lamberton, tried to establish a theocracy at the mouth of the Schuylkill River. The New Haven Colony had earlier struck a deal with the Lenape to buy much of New Jersey south of present-day Trenton. The Dutch and Swedes in the area burned the English colonists' buildings. A Swedish court under Swedish Governor Johan Björnsson Printz convicted Lamberton of "trespassing, conspiring with the Indians." The offshoot New Haven colony received no support. The Puritan Governor John Winthrop said it was dissolved owing to summer "sickness and mortality." The disaster contributed to New Haven's losing control of its area to the larger Connecticut Colony. [ citation needed ] In 1644, New Sweden supported the Susquehannock in their successful conflict with Maryland colonists (led by General Harrison II). The Dutch never recognized the legitimacy of the Swedish claim and, in the late summer of 1655, Director-General Peter Stuyvesant of New Amsterdam mustered a military expedition to the Delaware Valley to subdue the rogue colony. Although the colonists had to recognize the authority of New Netherland, the Dutch terms were tolerant. The Swedish and Finnish settlers continued to enjoy a much local autonomy, having their own militia, religion, court, and lands. This official status lasted until the English capture of New Netherland in October 1664, and continued unofficially until the area was included in William Penn's charter for Pennsylvania in 1682. By 1682, the area of modern Philadelphia was inhabited by about fifty Europeans, mostly subsistence farmers. In 1681, as part of a repayment of a debt, King Charles II granted William Penn a charter for what would become the Pennsylvania colony. Shortly after receiving the charter, Penn said he would lay out "a large Towne or Citty in the most Convenient place upon the Delaware River for health & Navigation." Penn wanted the city to live peacefully in the area, without a fortress or walls, so he bought the land from the Lenape. The legend is that Penn made a treaty of friendship with Lenape chief Tammany under an elm tree at Shackamaxon, in what became the city's Kensington District. Penn envisioned a city where all people regardless of religion could worship freely and live together. Being a Quaker, Penn had experienced religious persecution. He also planned that the city's streets would be set up in a grid, with the idea that the city would be more like the rural towns of England than its crowded cities. The homes would be spread far apart and surrounded by gardens and orchards. The city granted the first purchasers land along the Delaware River for their homes. It had access to the Delaware Bay and Atlantic Ocean, and became an important port in the Thirteen Colonies. He named the city Philadelphia (philos, "love" or "friendship", and adelphos, "brother") it was to have a commercial center for a market, state house, and other key buildings. Penn sent three commissioners to supervise the settlement and to set aside 10,000 acres (40 km 2 ) for the city. The commissioners bought land from Swedes at the settlement of Wicaco, and from there began to lay out the city toward the north. The area went about a mile along the Delaware River between modern South and Vine Streets. Penn's ship anchored off the coast of New Castle, Delaware, on October 27, 1682, and he arrived in Philadelphia a few days after that. He expanded the city west to the bank of the Schuylkill River, for a total of 1,200 acres (4.8 km 2 ). Streets were laid out in a gridiron system. Except for the two widest streets, High (now Market) and Broad, the streets were named after prominent landowners who owned adjacent lots. The streets were renamed in 1684 the ones running east–west were named after local trees (Vine, Sassafras, Mulberry, Cherry, Chestnut, Walnut, Locust, Spruce, Pine, Lombard, and Cedar) and the north–south streets were numbered. Within the area, four squares (now named Rittenhouse, Logan, Washington and Franklin) were set aside as parks open for everyone. Penn designed a central square at the intersection of Broad and what is now Market Street to be surrounded by public buildings. Some of the first settlers lived in caves dug out of the river bank, but the city grew with construction of homes, churches, and wharves. The new landowners did not share Penn's vision of a non-congested city. Most people bought land along the Delaware River instead of spreading westward towards the Schuylkill. The lots they bought were subdivided and resold with smaller streets constructed between them. Before 1704, few settlers lived west of Fourth Street. Philadelphia grew from a few hundred European inhabitants in 1683 to over 2,500 in 1701. The population was mostly English, Welsh, Irish, Germans, Swedes, Finns, and Dutch. Before William Penn left Philadelphia for the last time on October 25, 1701 he issued the Charter of 1701. The charter established Philadelphia as a city and gave the mayor, aldermen, and councilmen the authority to issue laws and ordinances and regulate markets and fairs. The first known Jewish resident of Philadelphia was Jonas Aaron, a German who moved to the city in 1703. He is mentioned in an article entitled "A Philadelphia Business Directory of 1703," by Charles H. Browning. It was published in The American Historical Register, in April, 1895. Philadelphia became an important trading center and major port. Initially the city's main source of trade was with the West Indies, which had established sugar cane plantations. It was part of the Triangle Trade, associated with Africa and Europe. During Queen Anne's War (1702 and 1713) with the French, trade was cut off to the West Indies, hurting Philadelphia financially. The end of the war brought brief prosperity to all of overseas British possessions, but a depression in the 1720s stunted Philadelphia's growth. The 1720s and '30s saw immigration from mostly Germany and northern Ireland to Philadelphia and the surrounding countryside. The region was developed for agriculture and Philadelphia exported grains, lumber products and flax seeds to Europe and elsewhere in the American colonies this pulled the city out of the depression. Philadelphia's pledge of religious tolerance attracted many other religions beside Quakers. Mennonites, Pietists, Anglicans, Catholics, and Jews moved to the city and soon outnumbered the Quakers, but they continued to be powerful economically and politically. Political tensions existed between and within the religious groups, which also had national connections. Riots in 1741 and 1742 took place over high bread prices and drunken sailors. In October 1742 and the "Bloody Election" riots, sailors attacked Quakers and pacifist Germans, whose peace politics were strained by the War of Jenkins' Ear. The city was plagued by pickpockets and other petty criminals. Working in the city government had such a poor reputation that fines were imposed on citizens who refused to serve an office after being chosen. One man fled Philadelphia to avoid serving as mayor. In the first half the 18th century, like other American cities, Philadelphia was dirty, with garbage and animals littering the streets. The roads were unpaved and in rainy seasons impassable. Early attempts to improve quality of life were ineffective as laws were poorly enforced. By the 1750s, Philadelphia was turning into a major city. Christ Church and the Pennsylvania State House, better known as Independence Hall, were built. Streets were paved and illuminated with oil lamps. Philadelphia's first newspaper, Andrew Bradford's American Weekly Mercury, began publishing on December 22, 1719. The city also developed culturally and scientifically. Schools, libraries and theaters were founded. James Logan arrived in Philadelphia in 1701 as a secretary for William Penn. He was the first to help establish Philadelphia as a place of culture and learning. Logan, who was the mayor of Philadelphia in the early 1720s, created one of the largest libraries in the colonies. He also helped guide other prominent Philadelphia residents, which included botanist John Bartram and Benjamin Franklin. Benjamin Franklin arrived in Philadelphia in October 1723 and would play a large part in the city's development. To help protect the city from fire, Franklin founded the Union Fire Company. In the 1750s Franklin was named one of the city's post master generals and he established postal routes between Philadelphia, New York, Boston, and elsewhere. He helped raise money to build the American colonies' first hospital, which opened in 1752. That same year the College of Philadelphia, another project of Franklin's, received its charter of incorporation. Threatened by French and Spanish privateers, Franklin and others set up a volunteer group for defense and built two batteries. When the French and Indian War began in 1754 as part of the Seven Years' War, Franklin recruited militias. During the war, the city attracted many refugees from the western frontier. When Pontiac's Rebellion occurred in 1763, refugees again fled into the city, including a group of Lenape hiding from other Native Americans, angry at their pacifism, and white frontiersmen. The Paxton Boys tried to follow them into Philadelphia for attacks, but was prevented by the city's militia and Franklin, who convinced them to leave. In the 1760s, the British Parliament's passage of the Stamp Act and the Townshend Acts, combined with other frustrations, increased political tension and anger against Britain in the colonies. Philadelphia residents joined boycotts of British goods. After the Tea Act in 1773, there were threats against anyone who would store tea and any ships that brought tea up the Delaware. After the Boston Tea Party, a shipment of tea had arrived in December, on the ship the Polly. A committee told the captain to depart without unloading his cargo. A series of acts in 1774 further angered the colonies activists called for a general congress and they agreed to meet in Philadelphia. The First Continental Congress was held in September in Carpenters' Hall. After the American Revolutionary War began in April 1775 following the Battles of Lexington and Concord, the Second Continental Congress met in May at the Pennsylvania State House. There they also met a year later to write and sign the Declaration of Independence in July 1776. Philadelphia was important to the war effort Robert Morris said, You will consider Philadelphia, from its centrical situation, the extent of its commerce, the number of its artificers, manufactures and other circumstances, to be to the United States what the heart is to the human body in circulating the blood. The port city was vulnerable to capture by the British by sea. Officials recruited soldiers and studied defenses for invasion from Delaware Bay, but built no forts or other installations. In March 1776 two British frigates began a blockade of the mouth of Delaware Bay British soldiers were moving south through New Jersey from New York. In December fear of invasion caused half the population to flee the city, including the Continental Congress, which moved to Baltimore. General George Washington pushed back the British advance at the battles of Princeton and Trenton, and the refugees and Congress returned. In September 1777, the British invaded Philadelphia from the south. Washington intercepted them at the Battle of Brandywine but was driven back. Thousands fled north into Pennsylvania and east into New Jersey Congress moved to Lancaster then to York. British troops marched into the half-empty Philadelphia on September 23 to cheering Loyalist crowds. The occupation lasted ten months. After the French entered the war on the side of the Continentals, the last British troops pulled out of Philadelphia on June 18, 1778, to help defend New York City. Continentals arrived the same day and reoccupied the city supervised by Major General Benedict Arnold, who had been appointed the city's military commander. The city government returned a week later, and the Continental Congress returned in early July. Historian Gary B. Nash emphasizes the role of the working class, and their distrust of their betters, in northern ports. He argues that working class artisans and skilled craftsmen made up a radical element in Philadelphia that took control of the city starting about 1770 and promoted a radical Democratic form of government during the revolution. They held power for a while, and used their control of the local militia to disseminate their ideology to the working class and to stay in power until the businessmen staged a conservative counterrevolution. Philadelphia suffered serious inflation, causing problems especially for the poor, who were unable to buy needed goods. This led to unrest in 1779, with people blaming the upper class and Loyalists. A riot in January by sailors striking for higher wages ended up with their attacking and dismantling ships. In the Fort Wilson Riot of October 4, men attacked James Wilson, a signer of the Declaration of Independence who was accused of being a Loyalist sympathizer. Soldiers broke up the riot, but five people died and 17 were injured. At the end of the American War of Independence, many Patriot soldiers had not been paid their wages for their service during the war. Congress refused the soldiers' request for payment of their salaries. In what is known as the Pennsylvania Mutiny of 1783, hundreds of Patriot veterans of the war who were owed back pay marched with their weapons on the Pennsylvania statehouse in Philadelphia. Congress, lacking in funds, fled from Philadelphia to Princeton, New Jersey. With their departure and the departure of their families and staffs, Philadelphia was left all but deserted. As a result of the Pennsylvania Mutiny of 1783, Congress fled Philadelphia, eventually settling in New York City, designated as the temporary capital. Besides the Constitutional Convention in May 1787, United States politics was no longer centered in Philadelphia. Due to political compromise, Congress chose a permanent capital to be built along the Potomac River. However, Philadelphia was selected as the temporary United States capital for ten years starting in 1790. The United States Congress, founded in March 1789, occupied the Philadelphia County Courthouse, which became known as Congress Hall, and the Supreme Court worked at City Hall. Robert Morris donated his home at 6th and Market Street as a residence for President Washington, known as the President's House. Yellow fever 1793 Edit After 1787, the city's economy grew rapidly in the postwar years. Serious yellow fever outbreaks in the 1790s interrupted development. Benjamin Rush identified an outbreak in August 1793 as a yellow fever epidemic, the first in 30 years, which lasted four months. Two thousand refugees from Saint-Domingue had recently arrived in the city in flight from the Haitian Revolution. They represented five percent of the city's total population. They likely carried the disease from the island where it was endemic, and it was rapidly transmitted by mosquito bites to other residents. Fear of contracting the disease caused 20,000 residents to flee the city by mid-September, and some neighboring towns prohibited their entry. Trade virtually stopped Baltimore and New York quarantined people and goods from Philadelphia. People feared entering the city or interacting with its residents. The fever finally abated at the end of October with the onset of colder weather and was declared at an end by mid-November. The death toll was 4,000 to 5,000, in a population of 50,000. Yellow fever outbreaks recurred in Philadelphia and other major ports through the nineteenth century, but none had as many fatalities as that of 1793. The 1798 epidemic in Philadelphia also prompted an exodus an estimated 1,292 residents died. Pennsylvania, which had abolished slavery in 1780, required any slaves brought to the city to be freed after six months' residency. The state law was challenged by French planters from Saint-Domingue, who brought their enslaved peoples with them, but defended by the Pennsylvania Abolition Society. Through 1796, 500 slaves from Saint-Domingue gained freedom in the city. Because of the violence accompanying the revolution on the island, Philadelphians, many of whom had southern ties, and residents of the Upper South worried that free people of color would encourage slave insurrections in the U.S. During the city's 10 years as federal capital, members of Congress were exempt from the abolition law, but the many slaveholders in the executive and judicial branches were not. President Washington, vice-President Jefferson and others brought slaves as domestic servants, and evaded the law by regularly shifting their slaves out of the city before the 6-month deadline. Two of Washington's slaves escaped from the President's House, and he gradually replaced his slaves with German immigrants who were indentured servants. The remains of the President's House were found during excavation for a new Liberty Bell Center, leading to archeological work in 2007. In 2010, a memorial on the site opened to commemorate Washington's slaves and African Americans in Philadelphia and U.S. history, as well as to mark the house site. The Pennsylvania state government left Philadelphia in 1799 and the United States government left in 1800. By this time, the city had become one of the United States' busiest ports and the country's largest city, with 67,787 people living in Philadelphia and its contiguous suburbs. Philadelphia's maritime trade was interrupted by the Embargo Act of 1807 and then the War of 1812. After the war, Philadelphia's shipping industry never returned to its pre-embargo status, and New York City succeeded it as the busiest port and largest city. The embargo and decrease in foreign trade led to the development of local factories to produce goods no longer available as imports. Manufacturing plants and foundries were built and Philadelphia became an important center of paper-related industries and the leather, shoe, and boot industries. Coal and iron mines, and the construction of new roads, canals, and railroads helped Philadelphia's manufacturing power grow, and the city became the United States' first major industrial city. Major industrial projects included the Waterworks, iron water pipes, a gasworks, and the U.S. Naval Yard. In response to exploitative working conditions, some 20,000 Philadelphia workers staged the first general strike in North America in 1835, in which workers in the city won the ten-hour workday and an increase in wages. In addition to its industrial power, Philadelphia was the financial center of the country. Along with chartered and private banks, the city was the home of the First and Second Banks of the United States, Mechanics National Bank and the first U.S. Mint. Cultural institutions, such as the Pennsylvania Academy of the Fine Arts, the Academy of Natural Sciences, the Athenaeum and the Franklin Institute also developed in the nineteenth century. The Pennsylvania General Assembly passed the Free School Law of 1834 to create the public school system. Ethnic rivalries Edit In the mid and late 1850s, immigrants from Ireland and Germany streamed into the city, swelling the population of Philadelphia and its suburbs. In Philadelphia, as the rich moved west of 7th Street, the poor moved into the upper class' former homes, which were converted into tenements and boarding houses. Many small row houses crowded alleyways and small streets, and these areas were filthy, filled with garbage and the smell of manure from animal pens. During the 1840s and 1850s, hundreds died each year in Philadelphia and the surrounding districts from diseases such as malaria, smallpox, tuberculosis, and cholera, related to poor sanitation the poor suffered the most fatalities. Small rowhouses and tenement housing were constructed south of South Street. Violence was a serious problem. Gangs like the Moyamensing Killers and the Blood Tubs controlled various neighborhoods. During the 1840s and early 1850s when volunteer fire companies, some of which were infiltrated by gangs, responded to a fire, fights with other fire companies often broke out. The lawlessness among fire companies virtually ended in 1853 and 1854 when the city took more control over their operations. During the 1840s and 50s violence was directed against immigrants by people who feared their competition for jobs and resented newcomers of different religions and ethnicity. Nativists often held mostly anti-Catholic, anti-Irish meetings. Violence against immigrants also occurred, the worst being the nativist riots in 1844. Violence against African Americans was also common during the 1830s, 40s, and 50s. Immigrants competed with them for jobs, and deadly race riots resulted in the burning of African-American homes and churches. In 1841, Joseph Sturge commented ". there is probably no city in the known world where dislike, amounting to the hatred of the coloured population, prevails more than in the city of brotherly love!" Several anti-slavery societies had been formed and free blacks, Quakers and other abolitionists operated safe houses associated with the Underground Railroad, but working class and ethnic whites opposed the abolitionist movement. The lawlessness and the difficulty in controlling it, along with residential development just north of Philadelphia, led to the Act of Consolidation in 1854. The act passed on February 2, making Philadelphia's borders coterminous with Philadelphia County, and incorporating various subdistrict within the county. Once the American Civil War began in 1861, Philadelphia's southern leanings were reduced. Popular hostility shifted against southern sympathizers. Mobs threatened a secessionist newspaper and the homes of suspected sympathizers, and were only turned away by the police and Mayor Alexander Henry. Philadelphia supported the war with soldiers, ammunition, and war ships and its manufacturers produced many army uniforms. Philadelphia was also a major receiving place of the wounded, with more than 157,000 soldiers and sailors treated within the city. Philadelphia began preparing for invasion in 1863, but the Confederate Army was repelled by Union forces at Gettysburg. In the years following the Civil War, Philadelphia's population continued to grow. The population grew from 565,529 in 1860 to 674,022 in 1870. By 1876, the city's population stood at 817,000. The dense population areas were not only growing north and south along the Delaware River, but also moving westward across the Schuylkill River. A large portion of the growth came from immigrants, still mostly Irish and German. In 1870, twenty-seven percent of Philadelphia's population was born outside the United States. In February 1854, the Act of Consolidation made the city of Philadelphia inclusive of the entire county, doing away with all other municipalities. By the 1880s, immigration from Russia, Eastern Europe, and Italy started rivaling immigration from Western Europe. Many of the immigrants from Russia and Eastern Europe were Jewish. In 1881, there were around 5,000 Jews in the city, and by 1905 the number had increased to 100,000. Philadelphia's Italian population grew from around 300 in 1870 to around 18,000 in 1900, with the majority settling in South Philadelphia. Along with foreign immigration, domestic migration by African Americans from the South led to Philadelphia having the largest black population of a Northern U.S. city in this period. By 1876, nearly 25,000 African Americans living in Philadelphia, and by 1890 the population was near 40,000. While immigrants moved into the city, Philadelphia's rich left for newer housing in the suburbs, with commuting made easy by newly constructed railroads. During the 1880s much of Philadelphia's upper class moved into the growing suburbs along the Pennsylvania Railroad's Main Line west of the city. Politically the city was dominated by the Republican Party, which had developed a strong political machine. The Republicans dominated the post-war elections, and corrupt officials made their way into the government and continued to control the city through voter fraud and intimidation. The Gas Trust was the hub of the city's political machine. The trust controlled the gas company supplying lighting to the city. With the board under complete control by Republicans in 1865, they awarded contracts and perks for themselves and their cronies. Some government reform did occur during this time. The police department was reorganized and volunteer fire companies were eliminated and were replaced by a paid fire department. A compulsory school act passed in 1895, and the Public School Reorganization Act freed the city's education from the political machine. Higher education changed as well. The University of Pennsylvania moved to West Philadelphia and reorganized to its modern form and Temple University, Drexel University and the Free Library were founded. The city's major project was organizing and staging the Centennial Exposition, the first World's Fair in the United States, which celebrated the nation's Centennial. Held in Fairmount Park, exhibits included Alexander Graham Bell's telephone and the Corliss Steam Engine. Beginning May 10, 1876, by the end of the Exposition on November 10, more than nine million people had visited the fair. The city undertook construction of a new city hall, designed to match its ambitions. The project was graft-ridden and it took twenty-three years to complete. Upon completion of its tower in 1894, City Hall was the tallest building in Philadelphia, a position it maintained until One Liberty Place surpassed it in 1986. Philadelphia's major industries of the era were the Baldwin Locomotive Works, William Cramp & Sons Ship and Engine Building Company, and the Pennsylvania Railroad. Westward expansion of the Pennsylvania Railroad helped Philadelphia keep up with nearby New York City in domestic commerce, as both cities fought for dominance in transporting iron and coal resources from Pennsylvania. Philadelphia's other local railroad was the Reading Railroad, but after a series of bankruptcies, it was taken over by New Yorkers. The Panic of 1873, which occurred when the New York City branch of the Philadelphia bank Jay Cooke and Company failed, and another panic in the 1890s hampered Philadelphia's economic growth. While the depressions hurt the city, its diverse array of industries helped it weather difficult times. It had numerous iron and steel-related manufacturers, including Philadelphian-owned iron and steel works outside the city, most notably the Bethlehem Iron Company in the city by that name. The largest industry in Philadelphia was textiles. Philadelphia produced more textiles than any other U.S. city in 1904 the textile industry employed more than 35 percent of the city's workers. The cigar, sugar, and oil industries also were strong in the city. During this time the major department stores: Wanamaker's, Gimbels, Strawbridge and Clothier, and Lit Brothers, were developed along Market Street. By the end of the century, the city provided nine municipal swimming pools, making it a leader in the nation. In the beginning of the 20th century Philadelphia had taken on a poor reputation. People both inside and outside of the city commented that Philadelphia and its citizens were dull and contented with its lack of change. Harper's Magazine commented that "The one thing unforgivable in Philadelphia is to be new, to be different from what has been." In his pioneering 1899 work of urban sociology The Philadelphia Negro W. E. B. Du Bois had written, "Few large cities have such a disreputable record for misgovernment as Philadelphia." Du Bois's study found, in addition to general mismanagement and neglect, severe racial disparities in employment, housing, health, education, and criminal justice. These disparities persisted for example, between 1910 and 1920 the proportion of black citizens of Philadelphia who developed tuberculosis was four to six times that of whites. Along with an image of "dullness" and of poor governance practices, Philadelphia was known for its political corruption. The Republican-controlled political machine, run by Israel Durham, permeated all parts of city government. One official estimated that US$5 million was wasted each year from graft in the city's infrastructure programs. The majority of residents were Republican, but voter fraud and bribery were still common. In 1905, the city enacted election reforms, such as personal voter registration and the establishing primaries for all city offices. But, residents became complacent, and the city's political bosses continued in control. After 1907, Boss Durham retired and his successor, James McNichol, never controlled much outside North Philadelphia. The Vare brothers, George, Edwin, and William, had created their own organization in South Philadelphia. With no central authority, Senator Boies Penrose took charge. In 1910, infighting between McNichol and the Vares contributed to the reform candidate, Rudolph Blankenburg, to be elected mayor. During his administration, he made numerous cost-cutting measures and improvements to city services, but he served only one term. The machine again gained control. The policies of Woodrow Wilson's administration reunited reformers with the city's Republican Party and World War I temporarily halted the reform movement. In 1917, the murder of George Eppley, a police officer defending City Council primary candidate James Carey, ignited the reformers again. They passed legislation to reduce the City Council from two houses to one, and provided council members an annual salary. With the deaths of McNichol in 1917 and Penrose in 1921, William Vare became the city's political boss. In the 1920s the public flouting of Prohibition laws, mob violence, and police involvement in illegal activities led Mayor W. Freeland Kendrick to appoint Brigadier General Smedley Butler of the U.S. Marine Corps as director of public safety. Butler cracked down on bars and speakeasies and tried to stop corruption within the police force, but demand for liquor and political pressure made the job difficult, and he had little success. After two years, Butler left in January 1926 and most of his police reforms were repealed. On August 1, 1928, Boss Vare suffered a stroke, and two weeks later a grand jury investigation into the city's mob violence and other crimes began. Numerous police officers were dismissed or arrested as a result of the investigation, but no permanent change resulted. Strong support among some residents for the Democratic presidential candidate Al Smith, who was Catholic, marked the city's turning away in the 20th century from the Republican Party. Philadelphia continued to grow with immigrants coming from Eastern Europe and Italy, as well as African American migrants from the South. Foreign immigration was briefly interrupted by World War I. The demand for labor for the city's factories, including the new U.S. Naval Yard at Hog Island, which constructed ships, trains, and other items needed in the war effort, helped attract blacks in the Great Migration. In September 1918, cases of the influenza pandemic were reported at the Naval Yard and began to spread. The disease became widespread following the Philadelphia Liberty Loans Parade, which was attended by more than 200,000 people. Mortality on some days was several hundred people and, by the time the pandemic began to subside in October, more than 12,000 people had died. The rising popularity of automobiles led to widening of roads and creation of Northeast (Roosevelt) Boulevard in 1914, the Benjamin Franklin Parkway in 1918, the changing of many existing streets to one-way streets in the early 1920s, and construction of the Delaware River (Benjamin Franklin) Bridge to New Jersey in 1926. Philadelphia began to modernize, steel and concrete skyscrapers were constructed, old buildings were wired for electricity, and the city's first commercial radio station was founded. In 1907, the city constructed the first subway. It hosted the Sesqui-Centennial Exposition in South Philadelphia, and in 1928 opened the Philadelphia Museum of Art. In the three years after the stock market crashed in 1929, 50 Philadelphia banks closed. Of those only two were large, Albert M. Greenfield's Bankers Trust Company and the Franklin Trust Company. Savings and loan associations also faced trouble, with mortgages of 19,000 properties being foreclosed in 1932 alone. By 1934, 1,600 of 3,400 savings and loan associations had shut down. From 1929 to 1933, regional manufacturing fell by 45 percent factory payrolls fell by 60 percent retail sales fell by 40 percent. Worst hit of all was construction, where payrolls dropped 84 percent. Unemployment peaked in 1933, when 11.5 percent of whites, 16.2 percent of African Americans, and 19.1 percent of foreign-born whites were out of work. Mayor J. Hampton Moore blamed people's economic woes, not on the worldwide Great Depression, but on laziness and wastefulness, and claimed there was no starvation in the city. Soon after, he fired 3,500 city workers, instituted pay cuts, forced unpaid vacation, and reduced the number of contracts the city awarded. This saved Philadelphia millions of dollars, and the efforts kept the city from defaulting on its debts, but were unpopular among the unemployed. The city relied on state money to fund relief efforts. Moore's successor S. Davis Wilson instituted numerous programs financed by Franklin D. Roosevelt's New Deal's Works Progress Administration, despite condemning the program during his mayoral campaign. At the peak of WPA-financed jobs in 1936, 40,000 Philadelphians were employed under the program. With encouragement from the state government and labor's founding of the Congress of Industrial Organizations (CIO), Philadelphia became a union city. Many trade unions discriminated against African Americans for years, and they were closed out of some labor advances. Workers' dissatisfaction with conditions led to numerous strikes in the textile unions, and the CIO organized labor in other industries, with more strikes taking place. During the 1930s, the Democratic Party began to grow in Philadelphia, influenced by the leadership of the Roosevelt administration during the Depression. A newly organized Independent Democratic Committee reached out to residents. In 1936, the Democratic National Convention was held in Philadelphia. The majority of voters in the city reelected the Democrat Franklin D. Roosevelt as president they also voted for Democratic Congressmen and state representatives. City government continued to be dominated by Republicans, but the politicians were elected by small margins. The beginning of World War II in Europe and the threat of the U.S. becoming involved generated new jobs in defense-related industries. After the U.S. became involved in the war in 1941, the city mobilized. Philadelphia consistently met war bond quotas and when the war ended in 1945, 183,850 residents were in the U.S. armed forces. With so many men serving in the military, there had been a labor shortage businesses and industries hired women and workers from outside the city. In 1944, the Philadelphia Transportation Company promoted African Americans to positions as motormen and conductors (from which they had previously been excluded) on public transportation vehicles. Resentful, other PTC workers protested and began a strike that nearly immobilized the city. President Roosevelt sent troops to replace the striking workers. After a federal ultimatum, the workers returned after six days. After World War II ended, Philadelphia had a serious housing shortage. Around half of the city's housing had been built in the 19th century, and many units lacked proper sanitary facilities, were overcrowded, and in poor condition. Competition for housing, as African Americans (many had come to the city in the Great Migration from the South) and Puerto Ricans moved into new neighborhoods, resulted in racial tension. The wealthier middle-class residents, often white, continued to move out to the suburbs in what became called white flight. The population peaked at more than two million residents in 1950 afterward the city's population declined while that of the neighboring suburban counties grew. Some residents moved out of the region altogether due to restructuring of industry and loss of tens of thousands of jobs in the city. Philadelphia lost five percent of its population in the 1950s, three percent in the 1960s and more than thirteen percent in the 1970s. Manufacturing and other major Philadelphia businesses, which had supported middle-class lives for the working class, were moving out of the area or shutting down in industrial restructuring, including major declines in railroads. The city encouraged development projects in University City in West Philadelphia and the area around Temple University in North Philadelphia, it removed the "Chinese Wall" elevated railway, and developed Market Street East around the transportation hub. Some gentrification occurred, with restoration of properties in historic neighborhoods such as Society Hill, Rittenhouse Square, Queen Village, and the Fairmount area. A non profit group Action Philadelphia was formed to improve and promote Philadelphia's image.The airport expanded, the Schuylkill Expressway and the Delaware Expressway (Interstate 95) were built, SEPTA was formed, and residential and industrial development took place in Northeast Philadelphia. Preparations for the United States Bicentennial in 1976 began in 1964. By the early 1970s, US$3 million had been spent but no plans were set. The planning group was reorganized and numerous citywide events were planned. Independence National Historical Park was restored and development of Penn's Landing was completed. Less than half the expected visitors came to the city for the Bicentennial, but the event helped revive the identity of the city, inspiring annual neighborhood events and fairs. In 1947, Richardson Dilworth was selected as the Democratic candidate, but lost to incumbent mayor Bernard Samuel. During the campaign Dilworth made numerous specific charges about corruption within city government. The City Council set up a committee to investigate, with findings followed by a grand jury investigation. The five-year investigation and its findings garnered national attention. US$40 million in city spending was found to be unaccounted for, and the president judge of the Court of Common pleas had been tampering with court cases. The fire marshal went to prison and an official in the tax collection office, a water department employee, a plumbing inspector, and head of the police vice squad each committed suicide after criminal exposures. The public and the press demanded reform and by the end of 1950, a new city charter was drafted. The new charter strengthened the position of the mayor and weakened the City Council. The council would be made of ten councilmen elected by district and seven at large. City administration was streamlined and new boards and commissions were created. In 1951, Joseph S. Clark was elected as the first Democratic mayor in 80 years. Clark filled administration positions based on merit and worked to weed out corruption. Despite reforms and the Clark administration, a powerful Democratic patronage organization eventually replaced the old Republican one. Clark was succeeded by Richardson Dilworth, who continued the policies of his predecessor. Dilworth resigned to run for governor in 1962, and city council president James H. J. Tate was elected as the city's first Irish Catholic mayor. Tate was elected mayor in 1963 and reelected in 1967 despite opposition from reformers who opposed him as an organization insider. As elsewhere in major US cities, the 1960s was a turbulent decade for the city. Numerous civil rights and anti-war protests took place, including large protests led by Marie Hicks to desegregate Girard College. Students took over the Community College of Philadelphia in a sit-in, race riots broke out in Holmesburg Prison, and a 1964 riot along West Columbia Avenue killed two people, injured over 300 and caused around US$3 million in damages. Crime was also a serious problem. Primarily drug-related gang warfare plagued the city, and in 1970 crime was rated the city's number one problem in a City Planning Commission survey. The court system was overtaxed and the tactics of the police department under Police Commissioner Frank Rizzo were controversial. Frank Rizzo was credited with preventing the level of violence seen in other cities at the time and was elected mayor in 1971. The outspoken Rizzo, who was reelected in 1975, was a divisive figure who had loyal supporters and passionate opponents. Police and fire departments and cultural institutions were well supported under Rizzo, but other city departments like the Free Library, the Department of Welfare and Recreation, the City Planning Commission and the Streets Department experienced large cuts. The radical group called MOVE formed in 1972, and tension soon developed with city officials. The first major clash occurred in 1978 at the group's Powelton Village headquarters, resulting in the death of a police officer. Nine MOVE members were convicted at trial and sentenced to prison. In 1985, a stand-off occurred at the group's new headquarters in West Philadelphia, whose residents were believed to be armed resisters. The police dropped a satchel bomb on the house from a helicopter it set off a fire that killed eleven MOVE members, including five children, and destroyed sixty-two neighboring houses. Survivors sued the city in civil court and won damages. Crime continued to be a problem in the 1980s. Deadly Mafia warfare plagued South Philadelphia, drug gangs and crack houses invaded the slums of the city, and the murder rate skyrocketed. William J. Green became mayor in 1980, and in 1984 W. Wilson Goode became Philadelphia's first African-American mayor. Development continued in areas in Old City and South Street, and large modern skyscrapers of glass and granite, designed by nationally known architects, were constructed in Center City. City employee labor contracts signed during the Rizzo administration helped set up a city financial crisis that Green and Goode were unable to prevent. The city was near bankruptcy at the end of the 1980s. A group of Hmong refugees had settled in Philadelphia after the end of the 1970s Laotian Civil War associated with the Vietnam War. They were attacked in discriminatory acts, and the city's Commission on Human Relations held hearings on the incidents. Anne Fadiman, author of The Spirit Catches You and You Fall Down, said that lower-class residents resented the Hmong receiving a $100,000 federal grant for employment assistance when they were also out of work they believed that American citizens should be getting assistance. : 192 Between 1982 and 1984, three quarters of the Hmong people who had settled in Philadelphia left for other cities in the United States to join relatives living elsewhere. : 195 Vietnamese and other immigrants from Asia have settled in the city, many near the Italian Market area. In addition, numerous Hispanic immigrants from Central and South America have entered the city, settling in North Philadelphia. William Penn & Quakers Arrive New Religious Sects — At the time when the English colonies in America were being settled many new ideas had risen in Europe on the subject of religion. The common people had begun to think very freely on this subject and a number of new sects were formed. Everywhere there were state religions, kept up by the governments, and by these the members of the new sects were often badly treated, but no treatment was severe enough to make them give up their beliefs. Many of them were put in prison—and the prisons of those times were horrible places, dens of filth and sickness—but despite this the new sects continued to grow. Those who suffered on earth were sure that they would be rewarded in heaven. George Fox and His Doctrine — Among the new sects was one founded in 1648 by a poor shoemaker named George Fox, and preached by him throughout England at such times as he was out of prison. Great numbers came to hear him and soon there were thousands of converts to his doctrine. He did not believe in fighting, or in taking oaths, or that one man was better than another, or in show and ceremony of any kind, or in paying to support the state religion. His followers would not take off their hats before any man, even before the king, or speak of any man as "you," for they thought this was a sign of pride. With them every man was "thou" or "thee." The Friends or Quakers — These people called themselves "Friends," or "Children of Light," for they held that all truth came to them through the "inner light," not through men's teachings. God spoke to their hearts, they said, and in so doing was their guide. They would tremble or quake when they felt that the inner light had come to them, and from this they were soon spoken of as "Quakers." This title was given them in derision, but it came to be that by which they were everywhere known. They are still Friends among themselves, but Quakers to the world at large. How the Quakers Were Treated — Of all the sects the Quakers were treated the worst. The prisons were crowded with them and hundreds of them died in these dreadful places. Most of them were poor they would not resist the officers of the law if a prison door were thrown open they would not walk out but they would not obey any law that interfered with their religion, or pay to help support the state religion, and the government found them a difficult people to deal with. It is well that you should know something about the history and opinions of the Quakers, for they are the people to whom we owe the State of Pennsylvania. William Penn — There were certain persons of importance among the Quakers, and chief among these was a man named William Penn. He was the son of an admiral in the British army, Sir William Penn, who had lent money to the king and had much power at the king's court. The young man was handsome, manly, and well educated, and like his father was a friend of the king also of his brother, the Duke of York, to whom King Charles had given all the land along the Hudson and the Delaware Rivers, in America. But young Penn was a man of strong mind. He had heard a Quaker preacher named Thomas Lee and was soon full of the new ideas, which he talked about at home and abroad. His father was so angry that he turned his Quaker son out of the house and the law officers soon put him in prison. But nothing could stop him he preached, he wrote, he was in and out of prison he taught Quakerism in Germany, and next to George Fox he was the leading Quaker in Europe. A Refuge in America — There was only one place to which the ill-treated members of the new sects could look for peace and safety. This was in America. Many years before, the Pilgrims and Puritans of England had found homes in New England, where there was no one to disturb them. Later on the Catholics had come for safety to Maryland. And now, William Penn began to look across the sea to find a place of refuge for his friends and fellow sufferers. Early Quakers in America — Some Quakers had already made their way to New England, but the Puritans would not have them there. Some they hanged and others they banished, and in this cruel way got rid of "the troublesome new-comers." Later on, a number came to New Jersey, where they soon became so numerous that Penn took part with other Quakers in the purchase of that province. Some of these settlers crossed the Delaware to its western side. Thus when Penn reached America he found Quakers in his new province. The Indian Country — The time was now close at hand for the Quaker settlement of Pennsylvania. Some of the New Jersey settlers wrote to William Penn and told him that "the Indian country on the west side of the Delaware is most beautiful to look upon, and only wanteth a wise people to render it, like ancient Canaan, 'the glory of the earth." Penn wanted a home for his Quaker brethren where they would be quite free to worship God in their own way. Here was the land waiting for him. It had as yet only a few hundred settlers, Swedes, Dutch, and English. It might be made a great Quaker commonwealth. Penn's Grant of Land — Admiral Penn was now dead and William had become the heir of his estate. The admiral had loaned King Charles II sixteen thousand pounds, a sum which the king, who spent all the money he could get, was not likely soon to pay back. In 1680 William Penn asked King Charles to grant him a tract of land in America in payment of this debt. This he found the king quite willing to do. It was an easy way to get out of debt by giving away land that belonged to the Indians. At the same time it would kelp him to get rid of those obstinate Quakers who kept his law officers so busy. So he readily gave Penn the land asked for, and by the 4th of March, 1681, the charter to the new province was drawn up and ready to be signed. Penn himself wrote much of it, partly copying from the charter by which Maryland was granted to Lord Baltimore. Extent and Name of the New Province — The king proposed to give Penn a tract of land between Maryland on the south and New York on the north extending northward from the 40th to the 43rd degree of latitude, and five degrees in longitude from the Delaware westward. But what was then thought to be the 40th parallel of latitude did not prove to be so, and this mistake made much trouble in later years, since disputes arose as to the border line between Pennsylvania and Maryland. This trouble began at once, but its story must be told later on. As for the name of the new colony, Penn proposed to call it New Wales. When this name was rejected he proposed Sylvania, or "Woodland." To this "Penn" was added by those who drew up the charter. The new proprietor did not like this it was too much like worldly pride for his Quaker ideas but the king would not strike it out, and so the name stood as Pennsylvania, or "Penn's Woodland." Markham Takes Possession — As may be imagined, the Quakers of England were greatly pleased by this transaction. The charter was barely signed before numbers of them prepared to cross the ocean to this new land of refuge. Penn at once sent out his cousin, Colonel William Markham, to take possession and act as his deputy. He reached the Delaware about July 1, 1681, landing at the Swedish village of Upland. There he visited some of the Indian chiefs and purchased from them a considerable tract of land, being part of what is now Bucks County. For this he gave the Indians a large variety of goods, such as wampum, guns, blankets, pipes, and many other things. The Indians were quite satisfied with this sale. They had plenty of land but little of these goods, and they were very willing to exchange part of their land for these useful articles. Philadelphia Laid Out — During that year three ships loaded with settlers came up the Delaware. Commissioners were also sent over to select a suitable place for the large town which Penn proposed to build. They were told to examine Upland, but they chose for the new town a place farther north, where the Delaware ran close to a high bank, and another river, called Schuylkill by the Dutch, ran into it. Here was to be the city named by Penn [as] Philadelphia, a word which means 'Brotherly Love.' As laid out, it was two miles long, from the Delaware to the Schuylkill, and one mile wide, from Vine Street to Cedar (now South) Street. As is well known, the city many years ago extended beyond these narrow limits. The Good Ship Welcome — On the 27th of October, 1682, the good ship Welcome, with William Penn and about seventy emigrants on board, came to anchor in front of New Castle, a settlement of the Dutch and Swedes in what is now Delaware. About one hundred passengers had set sail, but thirty had died of smallpox on the voyage and been buried at sea. Two days later Upland was reached. Penn is said to have changed the name of this place to Chester at the suggestion of his friend Pearson, who had come from Chester, England. Penn Goes to Philadelphia — William Penn was very anxious to see the place where his new city had been laid out, and the story is told that he went up the river from Upland in an open boat in early November. Many settlers were there already, and as he passed up by the city front he could see the cave dwellings which had been dug in the river bank. Here [temporary] excavations were made and over them were built roofs of split trees, branches, and twigs, the whole usually covered with sods. The chimneys were made of stones, clay, and river grass. In these cave dwellings lived many of the settlers in some small degree of comfort while their houses were being built, and in one of them, at the foot of Sassafras (now Race) Street, was born John Key, the first English child born in Pennsylvania. Penn made the child a present of a city lot. Penn inspected the site of his new town, still covered with woodland, with much pleasure. Its streets were so far laid out only on paper, but he could see how well nature had fitted the site for a great city. His plan was to have every house built in the middle of a large plot, "so that there may be grounds on either side for gardens or orchards or fields, that it may be a green country town that will never be burned and always wholesome."’ There is little trace of this fine plan in the city today. Most of the early houses were of wood, but some were built of stone and had balconies and porches. The scene was a very busy one as the new town grew in size, the women helping the men in their building work, even sawing wood and carrying mortar. Arrival of Settlers — During 1682 more than two thousand settlers arrived, most of them landing at Chester and Philadelphia. They had suffered on the long voyage, but they had brought much property with them from England—furniture, tools, building materials, and provisions—and were ready to begin housekeeping at once. There was plenty to eat, for fish, deer, turkeys, ducks and other wild fowl were supplied at low rates by the Indians, who got along very well with these quiet, peace-loving people. Penn and the Indians — As for the Indians, we may be sure they were eager to see the great William Penn, of whom much had been told them. He was quite as glad to see them, with their alert forms and dignified faces. He walked about with them, sat in their wigwams and ate of their roasted hominy. And when they began to show how they could jump, it is said that he surprised them by out-jumping the best of them. Penn was then less than forty years of age and no doubt very active and agile. The Treaty with the Indians — Not much can be said of the famous treaty with the Indians, though a picture of this has been made, with Penn in the centre and the Indians sitting all around. Very likely there was such a treaty, and it may have taken place under the elm tree at Kensington, where a treaty monument now stands. The elm tree blew down long ago and only the monument marks the spot. No record was kept of this famous treaty and we do not know just what took place. But many years afterwards some of the Indians said: "We shall never forget the counsel that William Penn gave us though we cannot write, like the English, yet we can keep in memory what was said in our councils." Not while Penn lived was a drop of Quaker blood shed by an Indian, and when he died his [Indian] admirers showed great grief at the loss of the great and good Onas," their best friend among the [English]. The Grant of Delaware — Penn was wise enough to see that it would be best to have his province extend to the ocean, and for this purpose the Duke of York gave him the territory now forming the State of Delaware. He had laid out three counties—Philadelphia, Bucks, and Chester—and there were three counties in Delaware which for twenty years formed part of Pennsylvania. Afterwards Delaware got a legislature of its own, but it remained under the governor of Pennsylvania until the Revolution. Over all his grand domain William Penn had almost princely control, his charter giving him much more power than King Charles had kept for himself. For this royal domain all he had to pay the king, aside from the sixteen thousand pounds of debt, was two beaver skins a year and one-fifth of all the gold and silver he should find. As these metals were not found, the beaver skins covered the whole rent. Penn, however, bought from the Indians all the land he used, and he gave the Swedes who had farms on the site of Philadelphia. As much good land elsewhere, for he was too honest to think that the King had any right to give away what did not belong to him, its true owners being the Indians. The First Assembly — Penn had called a meeting of representatives of the people, to assemble at Chester on December 6. They did not all come, for many of them were too busy building and farming, but about forty came together on the day fixed. To this Assembly Penn offered a code of laws which he had drawn up before leaving England. There was to be complete religious liberty, though non-believers in Christ could not vote or hold office. Only property holders could vote but this excluded only servants and vagrants, since all others had property. All persons were forbidden to sell strong liquors to the Indians. The death penalty was limited to those guilty of murder and treason. Dueling was prohibited and the drunkard could be fined. Such was the "Great Law." It had much else in it, but these were its leading features. It formed the basis of the government of Pennsylvania during the colonial period. It was great in giving the people full religious liberty, which did not then exist in Europe. It also cut down the penalty of death to murder and treason. At that time there were many small crimes in England for which people could be hanged, and the laws everywhere were very severe. In this way William Penn proved himself a liberal and far-seeing man. The Plan of Philadelphia — William Penn did much more than to make laws for his new province. He wished to have a fine and handsome city and laid out Philadelphia with streets crossing each other at right angles and much wider than the streets of the cities of England. Those that ran east and west were given the names of trees in the forest around, as Chestnut, Walnut, Pine, etc. Those running north and south were known by numbers. There were to be a High Street passing through the center from river to river, and a Broad Street through the center north and south. Each of these was to be one hundred feet wide. In the centre of the city, where these streets crossed, was to be a square of ten acres, and in each quarter of the city squares of eight acres. These squares still exist, except the central one, on which now stands Philadelphia's great City Hall. Growth of the City — As has been said, not many settlers had come to the Delaware in the fifty years before Penn's arrival. Afterwards they came in large numbers. In 1683 nearly one hundred houses were built in Philadelphia, and two years afterwards there were six hundred houses with about three thousand people. Many others settled in the country between the Falls of Trenton and Chester and Marcus Hook. In the latter place the first Friends’ meeting-house was built. Most of the country dwellers planted Indian corn the first spring and had good crops in the autumn. Penn was proud of the promising growth of his colony, which increased more rapidly than any other in America. Before he returned to England, in 1684, there were about five thousand persons in the new province. New Land Bought from the Indians — Immigration was so rapid that Penn soon saw the need of more land than that purchased by Markham, and he bought another large tract from the Indians. They were quite willing to dispose of part of their forest in exchange for the goods of the [Europeans], though they would have had no use for money. The story told about this purchase is a tradition and we cannot be sure of its truth. It is stated that in this (or perhaps some other) purchase the land bought was to go as far back as a man could walk in three days. Penn and his friends, with a number of Indians, set out from the mouth of Neshaminv Creek to make the walk, going along in an easy fashion, now and then sitting down to rest and eat their crackers and cheese, and for the Indians to smoke. At the end of a day and a half they had reached a large spruce tree near Baker Creek. The party by this time were tired, so Penn said he had land enough and would leave the remainder for a future day. It was a sad time for the poor Indians when that day came, as will be seen further on. The Letitia House — In the summer of 1683 Penn built a house to which he gave the name of his daughter Letitia, also giving her name to the street on which it stood. This house has been moved to a beautiful location in Fairmount Park, where it has hosts of visitors. He lived in this humble mansion part of the time and here held the sessions of his Council, which was both a lawmaking body and a court. Here, in February, 1684, the Council tried a woman on the charge of witchcraft, William Penn sitting as judge. The jury of eight Friends brought in the verdict: "Guilty of having the common fame of a witch, but not guilty in form and manner as she stands indicted." That was the only trial for witchcraft ever held in Pennsylvania. Education and Immigration — An important action of Penn and his Council was to establish a school in which the young people of the city might gain some degree of education, the master chosen being Enoch Flower, who for twenty years had been a teacher in England. New settlers were now coming rapidly, about fifty ships arriving in 1683. And these were by no means all Englishmen. Many Welsh came, most of them Friends, who settled through the country around. And there were many Germans also, some of whom founded the village of Germantown. Some of these were Friends, others belonged to German sects, though these were like the Friends in some of their religious views. Pennsylvania During the American Revolution Pennsylvania played an extremely important role in the American Revolution. The First and Second Continental Congresses were convened in Philadelphia. This is where the Declaration of Independence was written and signed. Numerous key battles and events of the war occurred in the colony, including the crossing of the Delaware River, the Battle of Brandywine, the Battle of Germantown, and the winter encampment at Valley Forge. The Articles of Confederation were also drafted in Pennsylvania, the document that formed the basis of the new Confederation that was created at the end of the Revolutionary War.
Expository writing or academic writing is a type of writing that is used to convey information correctly and effectively to those who do not have prior knowledge of the topic. Knowledge and skills for expository writing are essential for students' future career, as professional members in society are asked to write a number of documents and reports. In this course, the instructor will help students learn how to write grammatically correct and coherent, logical texts on a given topic in English in a small class setting. The instructor is a native English speaker with a lot of experience in research and writing papers and patents at academic institutions and in the pharmaceutical industry. The aim of this course is to help students learn how to write grammatically correct and coherent, logical texts on research topics or any other topics of interest. By the end of this course, students should be able to: 1. give descriptions of people, machinery and objects in English. 2. give explanations for observations made in experiments in English. 3. compare and contrast objects of interest in English. 4. write grammatically correct and coherent, logical paragraphs in English. Academic writing, Expository writing \|Specialist skills\|\|✔ Intercultural skills\|\|✔ Communication skills\|\|✔ Critical thinking skills\|\|✔ Practical and/or problem-solving skills\| The flow of each class is basically as follows: (1) Introduction to each topic, (2) class or group discussion to unpack ideas and expand ideas on a specific topic, (3) session of technical vocabulary and definitions, (4) impromptu writing, and (5) homework writing assignment. \|Course schedule\|\|Required learning\| \|Class 1\|\|Level Check Interview/ Orientation/ Class Rules\|\|Students should understand the outline and class rules for the course.\| \|Class 2\|\|Introduction to Academic Writing\|\|Students should be able to identify differences between academic and non-academic writing.\| \|Class 3\|\|How to Write an Introduction\|\|Students should be able to write an introduction to describe a new topic or concept.\| \|Class 4\|\|Definitions, Vocabulary and Word Origin\|\|Students should be able to write simple or extended definitions for technical terms in their study field.\| \|Class 5\|\|How to Write About a Technical Process\|\|Students should be able to write a description about a technical process.\| \|Class 6\|\|How to Write Descriptions of Devices or Machinery\|\|Students should be able to write a description about a technical device or machinery.\| \|Class 7\|\|How to Write Instructions for a Procedure\|\|Students should be able to write instructions for a procedure.\| \|Class 8\|\|How to Tabulate Experimental Data\|\|Students should be able to write experimentally determined data in tabulated form.\| \|Class 9\|\|How to Prepare Visual Aids of Tabulated Data\|\|Students should be able to prepare visual aids from tabulated data.\| \|Class 10\|\|Comparing and Contrasting\|\|Students should be able to compare and contrast two objects or concepts based on characteristics.\| \|Class 11\|\|How to Write about Results\|\|Students should be able to write about results for a poster or paper.\| \|Class 12\|\|How to Write a Discussion about Results\|\|Students should be able to write a discussion about results for a poster or paper.\| \|Class 13\|\|Writing Conclusion Statements\|\|Students should be able to write conclusion statements for a poster or paper.\| \|Class 14\|\|How to Write an Abstract\|\|Students should be able to write an abstract for a poster or paper.\| \|Class 15\|\|How to Write a Review or Critique\|\|Students should be able to write reviews or critiques for an article or presentations.\| Great Writing 5: From Great Essays to Research, 3rd ed. (Keith Folse & Tison Pugh), Cengage, ISBN: 978-1285194967 Specified by the instructor as necessary. Students' knowledge of writing skills and their ability to apply the skills to writing situations will be assessed. Midterm and final exams 60%, exercise problems 40%. English learners at any level are welcome.
Vector Algebra x 13.1. Basic Concepts A vector V in the plane or in space is an arrow: it is determined by its length, denoted j V and its direction. Two arrows represent the same vector if they have the same length and are parallel (see ﬁgure 13.1). We use vectors to represent entities which are described by magnitude and direction. For example,https://url.theworksheets.com/3kc4 82 DownloadsPreview and Download ! 14 Calculate the magnitude of any vector’s horizontal and vertical components. 15 Draw a vector’s horizontal and vertical components. 16 Use trig to calculate a vector’s direction. 17 Calculate a vectors direction as a degree measurement combined with compass directions. 18 Calculate a vector’s magnitude using trig or Pythagorean theorem.https://url.theworksheets.com/8iw 167 DownloadsPreview and Download ! VECTOR ALGEBRA 1 Introduction Vector algebra is necessary in order to learn vector calculus. We are deal-ing with vectors in three-dimensional space so they have three components. The number of spatial variables that functions and vector components can depend on is therefore also three. I assume that the reader is familiar with vector addition and subtractionhttps://url.theworksheets.com/50pz 58 DownloadsPreview and Download ! so long as the size and direction of the vector is maintained. The sum of a series of vectors is drawn as the arrow with its tail at the location of the tail of the first vector, and its head at the head of the last vector. A vector that is the sum of other vectors is often called the resultant. We don’t ever really subtract vectors.https://url.theworksheets.com/5hwt 63 DownloadsPreview and Download ! Chapter – 10 : Vector Algebra Worksheet ‐ 1 1. Find a unit vector in the direction of a ,⃗ = 3î - 2 ĵ + 6 k à . 2. Find the angle between the vectors a ,⃗= î + ĵ + k à and b ,⃗= î - ĵ +k à. 3. If l a ,⃗ l = √3 , l b ,⃗ l = 2 and a ,⃗.b ,⃗ = 3, Find angle between a ,⃗ and b ,⃗. 4. Write the value of (ı̂.https://url.theworksheets.com/6fcx 44 DownloadsPreview and Download ! 10. The magnitude of the vector product of the vector î + ĵ +k à with a unit vector along the sum of vectors 2î +4 ĵ -5k à and λî E2 ĵ 3 k à is equal to √2 . Find the value of λ . 11. Show that four points with position vectors 4î +5 ĵ+k à, - ĵ - k à , 3î +9 ĵ+4 k à and - 4î + 4 ĵ +4k àhttps://url.theworksheets.com/5hx9 61 DownloadsPreview and Download ! VECTOR ALGEBRA SCHEMATIC DIAGRAM Topic Concept Degree of importance Refrence NCERT Text Book Edition 2007 Vector algebra (i)Vector and scalars * Q2 pg428 (ii)Direction ratio and direction cosines * Q 12,13 pg 440 (iii)Unit vector * * Ex 6,8 Pg 436 (iv)Position vector of a point and collinear vectors * * Q 15 Pg 440 , Q 11Pg440 , Q 16https://url.theworksheets.com/6fcy 38 DownloadsPreview and Download ! Basic Vector Algebra in 1. Vector Equality: Two vectors and are equal if and only if and . 2. Vector Addition: The sum of the vectors and is defined by. 3. Scalar Multiplication: Suppose is a vector and . Then the scalar product of is defined by. Example Find the sum of the following vectors. 1. , 2. ,https://url.theworksheets.com/2x32 91 DownloadsPreview and Download ! Directions: Solve the following problems algebraically on a separate sheet of paper. 17. A hiker walks 4.5 km in one direction, then makes a 45˚ turn to the right and walks another 6.4 km. What is the magnitude of her displacement?https://url.theworksheets.com/4kgx 46 DownloadsPreview and Download ! Chapter – 10 : Vector Algebra Worksheet – 2 1. What is the cosine of the angle ,which the vector î + ĵ + k à makes with y-axis ? 2. If l a ,⃗ l = √3 , l b ,⃗ l = 2 and angle between a ,⃗ and b ,⃗ is 60 4 3. If a , ,⃗ L 2 áı ȷ̂ E3k à and b ,⃗ L 3ı̂ȷ̂ F2k à then find angle between a ,⃗ Eb ,⃗ and a ,⃗ Fb ,⃗ . 4.https://url.theworksheets.com/5hwv 55 DownloadsPreview and Download !
Help your third graders master key math skills, and provide differentiated math practice through fun, hands-on activities in these 12 centers. Each full-color center focuses on a skill from a Common Core State Standards math domain, such as numbers and counting, operations, measurement and data, and geometry. The easy-to-assemble centers are stored in pocket folders, making them easy to use at a table, desk, or quiet area in the classroom. After a teacher or classroom aide models how to use a center, students can complete the activity independently, in pairs, or with the teacher, who may use the task to informally assess a student’s understanding. Each of the 12 centers includes: Take It to Your Seat: Common Core Math Centers has been updated to address the new advanced standards and provides: - full-color center mats and task cards - an overview with lesson objectives - a student direction page that explains the center activity - a response form or written practice activity - a center checklist to record each student’s progress The 12 centers cover the following third grade math skills: Numbers and Counting - an increase in the variety and complexity of activities to practice each skill in alignment with Common Core State Standards. - leveled tasks within some centers so that tasks progress in difficulty. - an illustrated math concept or rule for every center to support visual learners and keep the focus on the targeted math concept or skill. - Describe and extend number patterns - Know place values in four-digit numbers - Round numbers to the nearest 10 or 100 - Understand a fraction as a part of a whole - Recognize fractions and equivalent fractions Measurement and Data - Build fluency with multiplication and division facts within 100 - Apply the inverse relationship between multiplication and division for facts within 100 - Identify the correct operation to create an equation - Tell time to the nearest minute - Use graphs to show and interpret data - Know how to find the perimeter and area of a shape - Define and recognize quadrilaterals
Spanish American wars of independence \|Spanish American wars of Independence\| \|Part of Latin American wars of independence\| Decisive events of the war: Cortes de Cádiz (1812) (top left); Congress of Cúcuta (1821) (bottom left); Crossing of the Andes (1817) (bottom right); Battle of Tampico (1829) (top right). The Spanish American wars of independence were the numerous wars against Spanish rule in Spanish America that took place during the early 19th century, after the French invasion of Spain during Europe's Napoleonic Wars. These conflicts started in 1809 with short-lived governing juntas established in Chuquisaca and Quito opposing the composition of the Supreme Central Junta of Seville. When the Central Junta fell to the French invasion, in 1810, numerous new juntas appeared across the Spanish domains in the Americas. The conflicts among these colonies and with Spain eventually resulted in a chain of newly independent countries stretching from Argentina and Chile in the south to Mexico in the north in the first third of the 19th century. Cuba and Puerto Rico remained under Spanish rule until the Spanish–American War in 1898. The new republics from the beginning abolished the casta system, the Inquisition and nobility, and slavery was ended in all of the new nations within a quarter century. Criollos (those of Spanish descent born in the New World) and mestizos (those of mixed Indian and Spanish blood) replaced Spanish-born appointees in most political offices. Criollos remained at the top of a social structure which retained some of its traditional features culturally, if not legally. For almost a century thereafter, conservatives and liberals fought to reverse or to deepen the social and political changes unleashed by those rebellions. These conflicts were fought as both wars of national liberation and civil wars, since on the one hand the goal of one group of belligerents was the independence of the Spanish colonies, and on the other the majority of combatants on both sides were Spanish Americans and indigenous people, not Spaniards. While some Spanish Americans believed that independence was necessary, most who initially supported the creation of the new governments saw them as a mean to preserve the region's autonomy from the French. Over the course of the next decade, the political instability in Spain and the absolutist restoration under Ferdinand VII convinced more and more Spanish Americans of the need to formally establish independence from the mother country. The events in Spanish America were related to the other wars of independence in Haiti and Brazil. Brazil's independence, in particular, shared a common starting point with Spanish America's, since both conflicts were triggered by Napoleon's invasion of the Iberian Peninsula, which forced the Portuguese royal family to resettle in Brazil in 1807. The process of Latin American independence took place in the general political and intellectual climate that emerged from the Age of Enlightenment and that influenced all of the Atlantic Revolutions, including the earlier revolutions in the United States and France. A more direct cause of the Spanish American wars of independence were the unique developments occurring within the Kingdom of Spain and its monarchy during this period. - 1 Historical background - 2 Creation of new governments in Spain and Americas, 1808-1810 - 3 First phase of the wars of independence, 1810–1814 - 4 Royalist ascendancy, 1814–1820 - 5 Independence consolidated, 1820–1825 - 6 Last royalist bastions, 1825–1833 - 7 Effects of independence - 8 Overview - 9 See also - 10 Footnotes - 11 Further reading Several factors set the stage for wars of independence. First the Bourbon Reforms of the mid-eighteenth century introduced changes to the relationship of Spanish Americans to the Crown. In an effort to better control the administration and economy of the overseas possessions the Crown reintroduced the practice of appointing outsiders, almost all peninsulars, to the various royal offices throughout the empire. This meant that Spanish Americans lost the gains they had made in holding local offices as a result of the sale of offices during the previous century and a half. In some areas—such as Cuba, Río de la Plata and New Spain—the reforms had positive effects, improving the local economy and the efficiency of the government. In other areas, the changes in crown's economic and administrative policies led to tensions with locals, which at times erupted into open revolts, such as the Revolt of the Comuneros in New Granada and the Rebellion of Túpac Amaru II in Peru. Neither of these two eighteenth-century developments—the loss of high offices to Criollos and the revolts—were the direct causes of the wars of independence, which took place decades later, but they were important elements of the political background in which the wars took place. Other factors included Enlightenment thinking and the examples of the Atlantic Revolutions. The Enlightenment spurred the desire for social and economic reform to spread throughout Latin America and the Iberian Peninsula. Ideas about free trade and physiocratic economics were raised by the Enlightenment in Spain. The political reforms implemented and the many constitutions written both in Spain and throughout the Spanish world during the wars of independence were influenced by these factors. Creation of new governments in Spain and Americas, 1808-1810 Collapse of the Bourbon dynasty The Peninsular War was the trigger for the wars of independence. The Peninsular War began an extended period of instability in the world-wide Spanish Monarchy which lasted until 1823. Napoleon's removal of the Bourbon dynasty from the Spanish throne precipitated a political crisis. Although the Spanish world almost uniformly rejected Napoleon's plan to give the crown to his brother, Joseph, there was no clear solution to the lack of a king. Following traditional Spanish political theories on the contractual nature of the monarchy (see Philosophy of Law of Francisco Suárez), the peninsular provinces responded to the crisis by establishing juntas. The move, however, led to more confusion, since there was no central authority and most juntas did not recognize the presumptuous claim of some juntas to represent the monarchy as a whole. The Junta of Seville, in particular, claimed authority over the overseas empire, because of the province's historic role as the exclusive entrepôt of the empire. Rebellion against Spanish Rule This impasse was resolved through negotiations between the juntas and the Council of Castile, which led to the creation of a "Supreme Central and Governmental Junta of Spain and the Indies" on September 25, 1808. It was agreed that the traditional kingdoms of the peninsula would send two representatives to this Central Junta, and that the overseas kingdoms would send one representative each. These "kingdoms" were defined as "the viceroyalties of New Spain, Peru, New Granada, and Buenos Aires, and the independent captaincies general of the island of Cuba, Puerto Rico, Guatemala, Chile, Province of Venezuela, and the Philippines." This scheme was criticized for providing unequal representation to the overseas territories; nevertheless, throughout the end of 1808 and early 1809, the provincial capitals elected candidates, whose names were forwarded to the capitals of the viceroyalties or captaincies general. Several important and large cities were left without direct representation in the Supreme Junta. In particular Quito and Chuquisaca, which saw themselves as the capitals of kingdoms, resented being subsumed in the larger "kingdom" of Peru. This unrest led to the establishment of juntas in these cities in 1809, which were eventually quashed by the authorities within the year. An unsuccessful attempt at establishing a junta in New Spain was also stopped. In order to establish a more legitimate government, the Supreme Junta called for the convening of an "extraordinary and general Cortes of the Spanish Nation." The election scheme for the Cortes, based on provinces and not kingdoms, was more equitable and provided more time to determine what would be considered an overseas province. The dissolution of the Supreme Junta on January 29, 1810, because of the reverses suffered after the Battle of Ocaña by the Spanish forces paid with Spanish American money, set off another wave of juntas being established in the Americas. French forces had taken over southern Spain and forced the Supreme Junta to seek refuge in the island-city of Cadiz. The Junta replaced itself with a smaller, five-man council, the Council of Regency of Spain and the Indies. Most Spanish Americans saw no reason to recognize a rump government that was under the threat of being captured by the French at any moment, and began to work for the creation of local juntas to preserve the region's independence from the French. Junta movements were successful in New Granada (Colombia), Venezuela, Chile and Río de la Plata (Argentina). Less successful, though serious movements, also occurred in Central America. Ultimately, Central America, along with most of New Spain, Quito (Ecuador), Peru, Upper Peru (Bolivia), the Caribbean and the Philippine Islands remained in control of royalists for the next decade and participated in the Spanish Cortes effort to establish a liberal government for the Spanish Monarchy. First phase of the wars of independence, 1810–1814 The creation of juntas in Spanish America set the stage for the fighting that would afflict the region for the next decade and a half. Political fault lines appeared, and were often the causes of military conflict. On the one hand the juntas challenged the authority of all royal officials, whether they recognized the Regency or not. On the other hand, royal officials and Spanish Americans who desired to keep the empire together were split between liberals, who supported the efforts of the Cortes, and conservatives (often called "absolutists" in the historiography), who did not want to see any innovations in government. Finally, although the juntas claimed to carry out their actions in the name of the deposed king, Ferdinand VII, their creation provided an opportunity for people who favored outright independence to publicly and safely promote their agenda. The proponents of independence called themselves patriots, a term which eventually was generally applied to them. The idea that independence was not the initial concern is evidenced by the fact that few areas declared independence in the years after 1810. The congresses of Venezuela and New Granada did so in 1811 and also Paraguay in same year (14 and 15 of May 1811). Some historians explain the reluctance to declare independence as a "mask of Ferdinand VII": that is, that patriot leaders felt that they needed to claim loyalty to the deposed monarch in order to prepare the masses for the radical change that full independence eventually would entail. Nevertheless, even areas such as Río de la Plata and Chile, which more or less maintained de facto independence from the peninsular authorities, did not declare independence until quite a few years later, in 1816 and 1818, respectively. Overall, despite achieving formal or de facto independence, many regions of Spanish America were marked by nearly continuous civil wars, which lasted well into the 1820s. In Mexico, where the junta movement had been stopped in its early stages by a coalition of Peninsular merchants and government officials, efforts to establish a government independent of the Regency or the French took the form of popular rebellion, under the leadership of Miguel Hidalgo. Hidalgo was captured and executed in 1811, but a resistance movement continued, which declared independence from Spain in 1813. In Central America, attempts at establishing juntas were also put down, but resulted in significantly less violence. The Caribbean islands, like the Philippines on the other side of the world, were relatively peaceful. Any plots to set up juntas were denounced to the authorities early enough to stop them before they gained widespread support. Underlying social tensions had a great impact on the nature of the fighting. Rural areas were pitted against urban centers, as grievances against the authorities found an outlet in the political conflict. This was the case with Hidalgo's peasant revolt, which was fueled as much by discontent over several years of bad harvests as with events in the Peninsular War. Hidalgo was originally part of a circle of liberal urbanites in Querétaro, who sought to establish a junta. After this conspiracy was discovered, Hidalgo turned to the rural people of the Mexican Bajío to build his army, and their interests soon overshadowed those of the urban intellectuals. A similar tension existed in Venezuela, where the Spanish immigrant José Tomás Boves was able to form a nearly invincible, though informal, royalist army out of the Llanero, mixed-race, plains people, by seeking to destroy the white landowning class. Boves and his followers often disregarded the command of Spanish officials and were not concerned with actually reestablishing the toppled royal government, choosing instead to keep real power among themselves. Finally in the backcountry of Upper Peru, the republiquetas kept the idea of independence alive by allying with disenfranchised members of rural society and Native groups, but were never able to take the major population centers. This period witnessed increasingly violent confrontations between Spaniards and Spanish Americans, but this tension was often related to class issues or fomented by patriot leaders to create a new sense of nationalism. After being incited to rid the country of the gachupines (a disparaging term for Peninsulares), Hidalgo's forces indiscriminately massacred hundreds of Criollos and Peninsulares who had taken refuge at the Alhóndiga de Granaditas in Guanajuato. In Venezuela during his Admirable Campaign, Simón Bolívar instituted a policy of a war to the death—in which and royalist Spanish Americans would be purposely spared but even neutral Peninsulares would be killed—in order to drive a wedge between the two groups. This policy laid the ground for the violent royalist reaction under Boves. Often though, royalism or patriotism simply provided a banner to organize the aggrieved, and the political causes could be discarded just as quickly as they were picked up. The Venezuelan Llaneros switched to the patriot banner once the elites and the urban centers became securely royalist after 1815, and it was the royal army in Mexico that ultimately brought about that nation's independence. Regional rivalry also played an important role in the wars. The disappearance of a central, imperial authority—and in some cases of even a local, viceregal authority (as in the cases of New Granada and Río de la Plata)—initiated a prolonged period of balkanization in many regions of Spanish America. It was not clear which political units which should replace the empire, and there were no new national identities to replace the traditional sense of being Spaniards. The original juntas of 1810 appealed first, to sense of being Spanish, which was counterposed to the French threat; second, to a general American identity, which was counterposed to the Peninsula lost to the French; and third, to a sense of belonging to the local province, the patria in Spanish. More often than not, juntas sought to maintain a province's independence from the capital of the former viceroyalty or captaincy general as much as from the Peninsula itself. Armed conflicts broke out between the provinces over the question of whether some provinces were to be subordinate to others as they had been under the crown. This phenomenon was particularly evident in New Granada and Río de la Plata. This rivalry also leads some regions to adopt the opposite political cause to that chosen by their rivals. Peru seems to have remained strongly royalist in large part because of its rivalry with Río de la Plata, to which it had lost control of Upper Peru when the later was elevated to a viceroyalty in 1776. The creation of juntas in Río de la Plata allowed Peru to regain formal control of Upper Peru for the duration of the wars. Royalist ascendancy, 1814–1820 \|Part of a series of articles on\| \|Part of the Politics series on\| By 1815 the general outlines of which areas were controlled by royalists and pro-independence forces were established and a general stalemate set in the war. In areas where royalists controlled the main population centers, most of the fighting by those seeking independence was done by isolated guerrilla bands. In New Spain, the two main guerrilla groups were led by Guadalupe Victoria in Puebla and Vicente Guerrero in Oaxaca. In northern South America, New Granadan and Venezuelan patriots, under leaders such as Francisco de Paula Santander, Simón Bolívar, Santiago Mariño, Manuel Piar and José Antonio Páez, carried out campaigns in the vast Orinoco River basin and along the Caribbean coast, often with material aid coming from Curaçao and Haiti. Also, as mentioned above, in Upper Peru, guerrilla bands controlled the isolated, rural parts of the country. During this period, royalist forces made advances into New Granada, which they controlled from 1815 to 1819, and into Chile, which they controlled from 1814 to 1817. Except for royalist areas in the northeast and south, the provinces of New Granada had maintained independence from Spain since 1810, unlike neighboring Venezuela, where royalists and pro-independence forces had exchanged control of the region several times. To pacify Venezuela and to retake New Granada, Spain organized in 1815 the largest armed force it ever sent to the New World, consisting of 10,500 troops and nearly sixty ships. (See, Spanish reconquest of New Granada.) Although this force was crucial in retaking a solidly pro-independence region like New Granada, its soldiers were eventually spread out throughout Venezuela, New Granada, Quito, and Peru and were lost to tropical diseases, diluting their impact on the war. More importantly, the majority of the royalist forces were composed, not of soldiers sent from the peninsula, but of Spanish Americans. Overall, Europeans formed only about a tenth of the royalist armies in Spanish America, and only about half of the expeditionary units, once they were deployed in the Americas. Since each European soldier casualty was replaced by a Spanish American soldier, over time, there were more and more Spanish American soldiers in the expeditionary units. For example Pablo Morillo, commander in chief of the expeditionary force sent to South America, reported that he had only 2,000 European soldiers under his command in 1820; in other words, only half the soldiers of his expeditionary force were European. It is estimated that in the Battle of Maipú only a quarter of the royalist forces were European soldiers, in the Battle of Carabobo about a fifth, and in the Battle of Ayacucho less than 1% was European. Restoration of Ferdinand VII In March 1814, following with the collapse of the First French Empire, Ferdinand VII was restored to the Spanish throne. This signified an important change, since most of the political and legal changes made on both sides of the Atlantic—the myriad of juntas, the Cortes in Spain and several of the congresses in the Americas, and many of the constitutions and new legal codes—had been made in his name. Before entering Spanish territory, Ferdinand made loose promises to the Cortes that he would uphold the Spanish Constitution. But once in Spain he realized that he had significant support from conservatives in the general population and the hierarchy of the Spanish Catholic Church; so, on May 4, he repudiated the Constitution and ordered the arrest of liberal leaders on May 10. Ferdinand justified his actions by stating that the Constitution and other changes had been made by a Cortes assembled in his absence and without his consent. He restored the former legal codes and political institutions and promised to convene a new Cortes under its traditional form (with separate chambers for the clergy and the nobility), a promise never fulfilled. News of the events arrived through Spanish America during the next three weeks to nine months, depending on time it took goods and people to travel from Spain. Ferdinand's actions constituted a definitive de facto break both with the autonomous governments, which had not yet declared formal independence, and with the effort of Spanish liberals to create a representative government that would fully include the overseas possessions. Such a government was seen as an alternative to independence by many in New Spain, Central America, the Caribbean, Quito, Peru, Upper Peru and Chile. Yet the news of the restoration of the "ancien régime" did not initiate a new wave of juntas, as had happened in 1809 and 1810, with the notable exception of the establishment of a junta in Cuzco demanding the implementation of the Spanish Constitution. Instead most Spanish Americans were moderates who decided to wait and see what would come out of the restoration of normalcy. In fact, in areas of New Spain, Central America and Quito, governors found it expedient to leave the elected constitutional ayuntamientos in place for several years in order to prevent conflict with the local society. Liberals on both sides of the Atlantic, nevertheless, continued to conspire to bring back a constitutional monarchy, ultimately succeeding in 1820. The most dramatic example of transatlantic collaboration is perhaps Francisco Javier Mina's expedition to Texas and northern Mexico in 1816 and 1817. Spanish Americans in royalist areas who were committed to independence had already joined the guerrilla movements. However, Ferdinand's actions did set areas outside of the control of the crown on the path to full independence. The governments of these regions, which had their origins in the juntas of 1810, and even moderates there, who had entertained a reconciliation with the crown, now saw the need to separate from Spain if they were to protect the reforms they had enacted. Towards the end of this period the pro-independence forces made two important advances. In the Southern Cone, a veteran of the Spanish army with experience in the Peninsular War, José de San Martín, became the governor of the Province of Cuyo. He used this position to begin organizing an army as early as 1814 in preparation for an invasion of Chile. This was an important change in strategy after three United Provinces campaigns had been defeated in Upper Peru. San Martín's army became the nucleus of the Army of the Andes, which received crucial political and material support in 1816 when Juan Martín de Pueyrredón became Supreme Director of the United Provinces. In January 1817, San Martín was finally ready to advance against the royalists in Chile. Ignoring an injunction from the congress of the Río de la Plata not to move against Chile, San Martín together with General Bernardo O'Higgins Riquelme, later Supreme Director of Chile, led the Army over the Andes in a move that turned the tables on the royalists. By February 10, San Martín had control of northern and central Chile, and a year later, after a war with no quarter, the south. With the aid of a fleet under the command of former British naval officer Thomas Cochrane, Chile was secured from royalist control and independence was declared that year. San Martín and his allies spent the next two years planning an invasion of Peru, which began in 1820. In northern South America, after several failed campaigns to take Caracas and other urban centers of Venezuela, Simón Bolívar devised a similar plan in 1819 to cross the Andes and liberate New Granada from the royalists. Like San Martín, Bolívar personally undertook the efforts to create an army to invade a neighboring country, collaborated with pro-independence exiles from that region, and lacked the approval of the Venezuelan congress. Unlike San Martín, however, Bolívar did not have a professionally trained army, but rather a quickly assembled mix of Llanero guerrillas, New Granadan exiles led by Santander and British recruits. From June to July 1819, using the rainy season as cover, Bolívar led his army across the flooded plains and over the cold, forbidding passes of the Andes, with heavy losses—a quarter of the British Legion perished, as well as many of his Llanero soldiers, who were not prepared for the nearly 4,000-meter altitudes—but the gamble paid off. By August Bolívar was in control of Bogotá and its treasury, and gained the support of many in New Granada, which still resented the harsh reconquest carried out under Morillo. Nevertheless Santander found it necessary to continue the policy of the "war to the death" and carried out the execution of thirty-eight royalist officers who had surrendered. With the resources of New Granada, Bolívar became the undisputed leader of the patriots in Venezuela and orchestrated the union of the two regions in a new state called Colombia (Gran Colombia). Independence consolidated, 1820–1825 To counter the advances the pro-independence forces had made in South America, Spain prepared a second, large, expeditionary force in 1819. This force, however, never left Spain. Instead, it became the means by which liberals were finally able to reinstate a constitutional regime. On January 1, 1820, Rafael Riego, commander of the Asturias Battalion, headed a rebellion among the troops, demanding the return of the 1812 Constitution. His troops marched through the cities of Andalusia with the hope of extending the uprising to the civilian population, but locals were mostly indifferent. An uprising, however, did occur in Galicia in northern Spain, and from there it quickly spread throughout the country. On March 7, the royal palace in Madrid was surrounded by soldiers under the command of General Francisco Ballesteros, and three days later, on March 10, the besieged Ferdinand VII, now a virtual prisoner, agreed to restore the Constitution. Riego's Revolt had two significant effects on the war in the Americas. Militarily, the large numbers of reinforcements, which were especially needed to retake New Granada and defend the Viceroyalty of Peru, would never arrive. Furthermore, as the royalists' situation became more desperate in region after region, the army experienced wholesale defections of units to the patriot side. Politically, the reinstitution of a liberal regime changed the terms under which the Spanish government sought to engage the insurgents. The new government naively assumed that the insurgents were fighting for Spanish liberalism and that the Spanish Constitution could still be the basis of reconciliation between the two sides. The government implemented the Constitution and held elections in the overseas provinces, just as in Spain. It also ordered military commanders to begin armistice negotiations with the insurgents with the promise that they could participate in the restored representative government. New Spain and Central America In effect, the Spanish Constitution of 1812 adopted by the Cortes de Cadiz served as the basis for independence in New Spain and Central America, since in both regions it was a coalition of conservative and liberal royalist leaders who led the establishment of new states. The restoration of the Spanish Constitution and representative government was enthusiastically welcomed in New Spain and Central America. Elections were held, local governments formed and deputies sent to the Cortes. Among liberals, however, there was fear that the new regime would not last; and conservatives and the Church worried that the new liberal government would expand its reforms and anti-clerical legislation. This climate of instability created the conditions for the two sides to forge an alliance. This alliance coalesced towards the end of 1820 behind Agustín de Iturbide, a colonel in the royal army, who at the time was assigned to destroy the guerrilla forces led by Vicente Guerrero. In January 1821, Iturbide began peace negotiations with Guerrero, suggesting they unite to establish an independent New Spain. The simple terms that Iturbide proposed became the basis of the Plan of Iguala: the independence of New Spain (now to be called the Mexican Empire) with Ferdinand VII or another Bourbon as emperor; the retention of the Catholic Church as the official state religion and the protection of its existing privileges; and the equality of all New Spaniards, whether immigrants or native-born. The following month the other important guerrilla leader, Guadalupe Victoria, joined the alliance, and March 1 Iturbide was proclaimed head of a new Army of the Three Guarantees. The representative of the new Spanish government, Superior Political Chief Juan O'Donojú, who replaced the previous viceroys, arrived in Veracruz on July 1; but he found that royalists the entire country except for Veracruz, Mexico City and Acapulco. Since at the time that O'Donojú had left Spain, the Cortes was considering greatly expanding the autonomy of the overseas Spanish possessions, O'Donojú proposed to negotiate a treaty with Iturbide on the terms of the Plan of Iguala. The resulting Treaty of Córdoba, which was signed on August 24, kept all existing laws, including the 1812 Constitution, in force until a new constitution for Mexico could be written. O'Donojú became part of the provisional governing junta until his death on October 8. Both the Spanish Cortes and Ferdinand VII rejected the Treaty of Córdoba, and the final break with the mother country came on May 19, 1822, when the Mexican Congress conferred the throne on Itrubide. Central America gained its independence along with New Spain. The regional elites supported the terms of the Plan of Iguala and orchestrated the union of Central America with the Mexican Empire in 1821. Two years later, following Iturbide's downfall, the region, with the exception of Chiapas, peacefully seceded from Mexico in July 1823, establishing the Federal Republic of Central America. The new state existed for seventeen years, centrifugal forces pulling the individual provinces apart by 1840. Unlike in New Spain and Central America, in South America independence was spurred by the pro-independence fighters who had held out for the past half decade. José de San Martín and Simón Bolívar inadvertently led a continent-wide pincer movement from southern and northern South America that liberated most of the Spanish American nations on that continent. After securing the independence of Chile in 1818, San Martín concentrated on building a naval fleet in the Pacific to counter Spanish control of those waters and reach the royalist stronghold of Lima. By mid-1820 San Martín had assembled a fleet of eight warships and sixteen transport ships under the command of Admiral Cochrane. The fleet set sail from Valparaíso to Paracas in southern Peru. On September 7, the army landed at Paracas and successfully took Pisco. After this, San Martín, waiting for a generalized Peruvian revolt, chose to avoid direct military confrontation. San Martín hoped that his presence would initiate an authentic Peruvian revolt against Spanish rule, believing that otherwise any liberation would be ephemeral. In the meantime, San Martín engaged in diplomacy with Viceroy Joaquín de la Pezuela, who was under orders from the constitutional government to negotiate on the basis of the 1812 Constitution and to maintain the unity of the Spanish Monarchy. However, these efforts proved fruitless, since independence and unity of the monarchy could not be reconciled, so the army sailed in late October to a better strategic position in Huacho, in northern Peru. During the next few months, successful land and naval campaigns against the royalists secured the new foothold, and it was at Huacho that San Martín learned that Guayaquil (in Ecuador) had declared independence on October 9. Bolívar, learning about the collapse of the Cadiz expedition, spent the year 1820 preparing a liberating campaign in Venezuela. Bolívar was aided by Spain's new policy of seeking engagement with the insurgents, which Morillo implemented, renouncing to the command in chief, and returning to Spain. Although Bolívar rejected the Spanish proposal that the patriots rejoin Spain under the Spanish Constitution, the two sides established a six-month truce and the regularization of the rules of engagement under the law of nations on November 25 and 26. The truce did not last six months. It was apparent to all that the royalist cause had been greatly weakened by the lack of reinforcements. Royalist soldiers and whole units began to desert or defect to the patriots in large numbers. On January 28, 1821, the ayuntamiento of Maracaibo, declared the province an independent republic that chose to join the new nation state of Gran Colombia. Miguel de la Torre, who had replaced Morillo as head of the army, took this to be a violation of the truce, and although the republicans argued that Maracaibo had switched sides of its own volition, both sides began to prepare for renewed war. The fate of Venezuela was sealed when Bolívar returned there in April leading an army of 7,000 from New Granada. At the Battle of Carabobo on June 24, the Gran Colombian forces decisively defeated the royalist forces, assuring control of Venezuela save for Puerto Cabello and guaranteeing Venezuelan independence. Bolívar could now concentrate on Gran Colombia's claims to southern New Granada and Quito. In Peru, on January 29, 1821, Viceroy Pezuela was deposed in a coup d'état by José de la Serna, but it would be two months before San Martín moved his army closer to Lima by sailing it to Ancón. During the next few months San Martín once again engaged in negotiations, offering the creation of an independent monarchy; but La Serna insisted on the unity of the Spanish monarchy, so the negotiations came to nothing. By July La Serna judged his hold on Lima to be weak, and on July 8 the royal army abandoned the coastal city in order to reinforce positions in the highlands, with Cuzco as new capital of viceroyalty. On the 12th San Martín entered Lima, where he was declared "Protector of the Country" on July 28, an office which allowed him to rule the newly independent state. To ensure that the Presidency of Quito became a part of Gran Colombia and did not remain a collection of small, divided republics, Bolívar sent aid in the form of supplies and an army under Antonio José de Sucre to Guayaquil in February 1821. For a year Sucre was unable to take Quito, and by November both sides, exhausted, signed a ninety-day armistice. The following year, at Battle of Pichincha on May 24, 1822, Sucre's Venezuelan forces finally conquered Quito; Gran Colombia's hold on the territory was secure. The following year, after a Peruvian patriot army was destroyed in the Battle of Ica, San Martín met with Simón Bolívar in Guayaquil on July 26 and 27. Thereafter San Martín decided to retire from the scene. For the next two years, two armies of Rioplatense (Argentinian), Chilean, Colombian and Peruvian patriots were destroyed trying to penetrate the royalist bastion in the Andean regions of Peru and Upper Peru. A year later a Peruvian congress resolved to make Bolívar head of the patriot forces in the country. An internecine conflict between La Serna and General Pedro Antonio Olañeta, which was an extension of the Liberal Triennium, proved to be the royalists' undoing. La Serna lost control of half of his best army by the beginning of 1824, giving the patriots an opportunity. Under the command of Bolivar and Sucre, the experienced veterans of the combined army, mainly Colombians, destroyed a royalist army under La Serna's command in the Battle of Ayacucho on December 9, 1824. La Serna's army was numerically superior but consisted of mostly new recruits. The only significant royalist area remaining on the continent was the highland country of Upper Peru. Following the Battle of Ayacucho, the royalist troops of Upper Peru under the command of Olañeta surrendered after he died in Tumusla on April 2, 1825. Bolívar tended to favor maintaining the unity of Upper Peru with Peru, but the Upper Peruvian leaders—many former royalists, like Casimiro Olañeta, nephew of General Olañeta—gathered in a congress under Sucre's auspices supported the country's independence. Bolívar left the decision to Sucre, who went along with the congress. Sucre proclaimed Upper Peru's independence in the city which now bears his name on August 6, bringing the main wars of independence to an end. As it became clear that there was to be no reversal of Spanish American independence, several of the new states began to receive international recognition. Early, in 1822, the United States recognized Chile, the United Provinces of the Río de la Plata, Peru, Gran Colombia, and Mexico. Britain waited until 1825, after the Battle of Ayacucho, to recognize Mexico, Gran Colombia, and Río de la Plata. Both nations recognized more Spanish American states in the next few years. Last royalist bastions, 1825–1833 The Spanish coastal fortifications in Veracruz, Callao and Chiloé, were the footholds that resisted until 1825 and 1826 respectively. In the following decade, royalist guerrillas continued to operate in several countries and Spain launched a few attempts to retake parts of the Spanish American mainland. In 1827 Colonel José Arizabalo started an irregular war with Venezuelan guerrillas, and Brigadier Isidro Barradas lead the last attempt with regular troops to reconquer Mexico in 1829. The Pincheira brothers moved to Patagonia and remained there as royalist outlaws until defeated in 1832. But efforts like these did not reverse the new political situation. The increasing irrelevancy of the Holy Alliance after 1825 and the fall of absolutism in France in 1830 during the July Revolution eliminated the principal support of Ferdinand VII in Europe, but it was not until the king's death in 1833 that Spain finally abandoned all plans of military re-conquest, and in 1836 its government went so far as to renounce sovereignty over all of continental America. During the course of the 19th century, Spain would recognize each of the new states. Only Cuba and Puerto Rico remained under Spanish rule, until the Spanish–American War in 1898. Effects of independence The nearly decade and a half of wars greatly weakened the Spanish American economies and political institutions, which hindered the region's potential economic development for most of the nineteenth century and resulted in the enduring instability the region experienced. Independence destroyed the de facto trade bloc that was the Spanish Empire - Manila galleons and Spanish treasure fleets in particular. After independence, trade among the new Spanish American nations was less than it had been in the colonial period. Once the ties were broken, the small populations of most of the new nations provided little incentive to entice Spanish American producers to recreate the old trade patterns. In addition, the protection against European competition, which the Spanish monopoly had provided to the manufacturing sectors of the economy, ended. Due to expediency, protective tariffs for these sectors, in particular textile production, were permanently dropped and foreign imports beat out local production. This greatly affected Native communities, which in many parts of Spanish America, specialized in supplying finished products to the urban markets, albeit using pre-industrial techniques. The wars also greatly affected the principal economic sector of the region, mining. Silver production in Bolivia halved after independence and it dropped by three quarters in Mexico. To compensate for the lack of capital, foreign investment, in particular from Great Britain, was courted, but it was not sizable enough to initiate an economic recovery. Finally the new nations entered the world economy after the end of the French Revolutionary and Napoleonic Wars, when the economies of Europe and the United States were recovering and aggressively seeking new markets to sell their products after more than two decades of disruption. Ultimately Spanish America could only connect to the world markets as an exporter of raw materials and a consumer of finished products. In addition to improving the economy, the lower social classes also had to be integrated into the new body politic, although they often got few rewards from independence. The political debate seeking answers to these questions was marked by a clash—at times on the battlefield—between liberalism and conservatism. Conservatives sought to maintain the traditional social structures in order to ensure stability; liberals sought to create a more dynamic society and economy by ending ethnically-based social distinctions and freeing property from economic restrictions. In its quest to transform society, liberals often adopted policies that were not welcome by Native communities, who had benefited from unique protections afforded to them by traditional Spanish law. Independence, however, did initiate the abolition of slavery in Spanish America, as it was seen as part of the independence struggle, since many slaves had gained their manumission by joining the patriot armies. In areas where slavery was not a major source of labor (Mexico, Central America, Chile), emancipation occurred almost immediately after independence was achieved. In areas where slavery was a main labor source(Colombia, Venezuela, Peru, Argentina), emancipation was carried out in steps over the next three decades, usually first with the creation of free-womb laws and programs for compensated emancipation. By the early 1850s, slavery had been abolished in the independent nations of Spanish America. Role of women Women were not simply spectators throughout the Independence Wars of Latin America. Many women took sides on political issues and joined independence movements in order to participate on many different levels. Women could not help but act as caring relatives either as mother, sister, wives or daughters of the men who were fighting. Women created political organizations and organized meetings and groups to donate food and supplies to the soldiers. Some women supported the wars as spies, informants and combatants. Manuela Sáenz was a long term lover of Simón Bolívar and acted as his spy and confidante and was secretary of his archive. She saved his life on two occasions, nursed wounded soldiers and has even been believed some historians to have fought in a few battles. Sáenz followed Bolivar and his army through the independence wars and became to be known in Latin America as the “mother of feminism and women’s emancipation and equal rights.” Bolivar himself was a supporter of women’s rights and suffrage in Latin America. It was Bolivar who allowed for Sáenz to become the great pioneer of women’s freedom. He wanted to set the women of Latin America free from the oppression and inferiority of what the Spanish regime had established. Bolivar even made Sáenz a Colonel of the Colombian Army due to her heroics which caused controversy because there were no women in the army at the time. Another woman who gained prominence in the fight for independence was Juana Azurduy de Padilla, a mixed-race woman who fought for the independence of in the Rio de la Plata region. Argentine President Cristina Fernández de Kirchner posthumously promoted to the rank of general. According to gender stereotypes, women were not meant to be soldiers; only men were supposed to engage in the fighting and conflict. There were still plenty of women presence on the battlefields to help rescue and nurse soldiers. Some women fought alongside their husbands and sons on the battlefield. The majority of women assumed supportive and non-competitive roles such as fund raising and caring for the sick. Revolution for women meant something differently than to men. Women saw revolution as a way to earn equal rights as men, such as voting, and to overcome the suppression of the superiority of men over women. Women were usually identified as victims during the independence wars for the women of Latin America were forced to sacrifice for the cause. The ideals of womanhood meant that women must sacrifice what the situation required such as a mother sacrificing her son or a virgin knowing she might be sacrificing motherhood or being wife due to the loss of many young men. This view meant that women were meant to contribute to independence in a supportive role while leaving the combat and politics in the hands of the men. Government and politics Independence also did not result in stable political regimes, save in a few countries. First, the new nations did not have well-defined identities, but rather the process of creating identities was only beginning. This would be carried out through newspapers and the creation of national symbols, including new names for the countries ("Mexico", "Colombia," "Ecuador," "Bolivia," "Argentina"), that broke with the past. In addition, the borders were also not firmly established, and the struggle between federalism and centralism, which begun in independence, continued throughout the rest of the century. Two large states that emerged from the wars—Gran Colombia and the Federal Republic of Central America—collapsed after a decade or two, and Argentina would not consolidate politically until the 1860s. The wars destroyed the old civilian bureaucracy that had governed the region for centuries, as institutions such as the audiencias were eliminated and many Peninsular officials fled to Spain. The Catholic Church, which had been an important social and political institution during the colonial period, initially came out weakened by the end of the conflicts. As with government officials, many Peninsular bishops abandoned their dioceses and their posts were not filled for decades until new prelates could be created and relations between the new nations and the Vatican was regularized. Then as the Church recovered, its economic and political power was attacked by liberals. Despite the fact that the period of the wars of independence itself was marked by a rapid expansion of representative government, for several of the new nations the nineteenth century was marked by militarism because of the lack of well-defined political and national institutions. The armies and officers that came into existence during the process of independence wanted to ensure that they got their rewards once the struggle was over. Many of these armies did not fully disband once the wars were over and they proved to be one of the stabler institutions in the first decades of national existence. These armies and their leaders effectively influenced the course of political development. Out of this new tradition came the caudillos, strongmen who amassed formal and informal economic, military and political power in themselves. Wars, battles and revolts \|New Spain and Guatemala\|\|New Granada, Venezuela, and Quito\| \|Río de la Plata, Paraguay and Upper Peru\|\|Chile and Peru\| \|New Spain, Guatemala, Cuba & Puerto Rico\|\|New Granada, Venezuela & Quito\|\|Río de la Plata, Montevideo & Paraguay\|\|Chile, Peru & Upper Peru\| - British Legions - Spanish reconquest of Mexico - Latin American wars of independence - Wars of national liberation - History of South America - Timeline of the Spanish American wars of independence - Garret, David T (2003). "Los incas borbónicos: la elite indígena cuzqueña en vísperas de Tupac Amaru". Revista Andina 36. ISSN 0259-9600. See also: - First invasion of Cisplatina by the Portuguese army lead by Diogo de Sousa on 1811 to annex the Banda Oriental to themselves, a colonial territory disputed between Spain and Portugal. Not for destroy the independent government of Buenos Aires. - Military units of Venezuela and Colombia with Irish and British volunteers or mercenaries under Latin American flags. - Lynch, The Spanish American Revolutions, 17–19, 334–335. Rodríguez, The Independence of Spanish America, 19–27. Kinsbruner, Independence in Spanish America, 7–12. - Lynch, Spanish American Revolutions, 5–17. Rodríguez, Independence of Spanish America, 24–25. Kinsbruner, Independence in Spanish America, 12–14, 17–32. - Lynch, Spanish American Revolutions, 27–34. Rodríguez, Independence of Spanish America, 14–18. Kinsbruner, Independence in Spanish America, 14–17, 23. - William Spence Robertson, "The Juntas of 1808 and the Spanish Colonies," English Historical Review (1916) 31#124 pp. 573-585 in JSTOR - Lynch, Spanish American Revolutions, 36–37. Rodríguez, Independence of Spanish America, 51–56, 58–59. Kinsbruner, Independence in Spanish America, 12, 35–37. - Royal Order of the Central Junta of January 22, 1809, cited in Rodríguez, Independence of Spanish America, 60. - Lynch, Spanish American Revolutions, 50–52, 236–239. Rodríguez, Independence of Spanish America, 53–55, 61–70, 80–81. Kinsbruner, Independence in Spanish America, 43–45. - "Batalla de Ocaña". Bicentenario de las independencias iberoaméricanas. Ministerio de Educación, Cultura y Deporte (Spain). Retrieved 2012-08-17. - Lynch, Spanish American Revolutions, 43–45, 52–56, 132–133, 195–196, 239–240. Rodríguez, Independence of Spanish America, 75–82, 110–112, 123–125, 136–139, 150–153. Kinsbruner, Independence in Spanish America, 36–37, 46, 52–53, 58–59, 61–62. - Lynch, Spanish American Revolutions, 36–37, 134–135. Rodríguez, Independence of Spanish America, 52–53. Kinsbruner, Independence in Spanish America, 45–46, 53. - The phrase is used by Lynch, Spanish American Revolutions, 56–58, 133. For a similar analysis without the phrase, see Crow, John A (1946). The Epic of Latin America. Garden City, N.Y.: Doubleday. pp. 425–426. - Lynch, Spanish American Revolutions, 107–111, 134–137, 162–172, 195–200, 238–240, 313–319, 335. Rodríguez, Independence of Spanish America, 93–111, 115, 123–126, 136–144, 147–156, 164–165, 168, 176–177. Kinsbruner, Independence in Spanish America, 46, 50, 52–53, 66–67, 100–101. - Lynch, Spanish American Revolutions, 118–121, 197–198, 200, 204–207, 306–313. Rodríguez, Independence of Spanish America, 113–122, 132, 159–167. Kinsbruner, Independence in Spanish America, 54, 66–70. - Lynch, Spanish American Revolutions, 121, 131–132. Rodríguez, Independence of Spanish America, 13–19, 22, - Lynch, Spanish American Revolutions, 57–71, 162–163, 240–242. Rodríguez, Independence of Spanish America, 111–113, 126–136, 153–159, 176–179. Kinsbruner, Independence in Spanish America, 53, 59. - Rodríguez, Independence of Spanish America, 168, 184, Kinsbruner, Independence in Spanish America, 70, 97. - Lynch, Spanish American Revolutions, 209. Rodríguez, Independence of Spanish America, 122. Kinsbruner, Independence in Spanish America, 57. - Small contingents from Spain had been arriving in the Americas since 1810. On August 25, 1810, a group of Spanish Marines arrived in Veracruz from Cádiz on the frigate, Nuestra señora de Atocha under the command of Rosendo Porlier and accompanying Viceroy Francisco Javier Venegas. These were the first Spaniards to have come from Europe in support of royalists. Frieyro de Lara. Guerra ejército y sociedad en el nacimiento de la España contemporánea. (2009, Universidad de Granada) p. 660. - Rebecca Earle, "'A Grave for Europeans'? Disease, Death, and the Spanish-American Revolutions" in Christon I. Archer, ed. The Wars of Independence in Spanish America, 283–297. - Rodríguez, Independence of Spanish America, 169–172. Kinsbruner, Independence in Spanish America, 56–57. - Lynch, Spanish American Revolutions, 336. Rodríguez, Independence of Spanish America, 106. - Lynch, Spanish American Revolutions, 162. 171–172, 207. Rodríguez, Independence of Spanish America, 173–175, 192–194 - Lynch, Spanish American Revolutions, 138–141. Rodríguez, Independence of Spanish America, 179–182. Kinsbruner, Independence in Spanish America, 72–75. - Lynch, Spanish American Revolutions, 209–218. MacKenzie, S. P. (1997). Revolutionary Armies in the Modern Era: A Revisionist Approach. London: Routledge. pp. 54, 61–64. ISBN 0-415-09690-1. Rodríguez, Independence of Spanish America, 184–192. Kinsbruner, Independence in Spanish America, 78–87. - Rodríguez, Independence of Spanish America, 194. Kinsbruner, Independence in Spanish America, 88, 114, 120–121, 127–128. - Lynch, Spanish American Revolutions, 335–340. Rodríguez, Independence of Spanish America, 194–195. Kinsbruner, Independence in Spanish America, 89. - Lynch analyzes the events through the older theory of a "conservative revolution": Spanish American Revolutions, 319–320. Compare to Rodríguez, Independence of Spanish America, 196–197, 199–205, 241–242. Kinsbruner, Independence in Spanish America, 97–98. Peter F. Guardino, "The War of Independence in Guerrero, New Spain, 1808–1821" in Archer, The Wars of Independence in Spanish America, 122–124. - Lynch, Spanish American Revolutions, 320–323. Rodríguez, Independence of Spanish America, 206–210. Kinsbruner, Independence in Spanish America, 98–99. Guardino, "The War of Independence in Guerrero," 121, 124–125. - Lynch, Spanish American Revolutions, 333–340. Rodríguez, Independence of Spanish America, 210–213. Kinsbruner, Independence in Spanish America, 100, 146–149. - Lynch, Spanish American Revolutions, 172–178. Rodríguez, Independence of Spanish America, 213–214. Kinsbruner, Independence in Spanish America, 76. - Lynch, Spanish American Revolutions, 218–219. Rodríguez, Independence of Spanish America, 219. Kinsbruner, Independence in Spanish America, 88–90. - Lynch, Spanish American Revolutions, 178–179. Rodríguez, Independence of Spanish America, 214–219. Kinsbruner, Independence in Spanish America, 76–77. - Lynch, Spanish American Revolutions, 185–189, 247–249, 267–272. Rodríguez, Independence of Spanish America, 219–220, 222–231. Timothy E. Anna, "Chaos and the Military Solution: The Fall of Royalist Government in Peru" in Archer, The Wars of Independence in Spanish America, 272–273. Kinsbruner, Independence in Spanish America, 77–78, 90–95. - Bushnell, David (1970). The Santander Regime in Gran Colombia. Westport: Greenwood Press. pp. 325–335. ISBN 0-8371-2981-8. Lynch, Spanish American Revolutions, 272–273, 279–284. Rodríguez, Independence of Spanish America, 232–234. Kinsbruner, Independence in Spanish America, 95–96. Chasteen, John Charles (2008). Americanos: Latin America's Struggle for Independence. Oxford University Press. pp. 164–165. ISBN 978-0-19-517881-4. - Kinsbruner, Independence in Spanish America, 105–106. - Costeloe, Michael P. Response to Revolution, 100 - Lynch, Spanish American Revolutions, 344–347. Rodríguez, Independence of Spanish America, 245. Kinsbruner, Independence in Spanish America, 131–136. - Lynch, Spanish American Revolutions, 343–344. Rodríguez, Independence of Spanish America, 244–245. Kinsbruner, Independence in Spanish America, 133–136. - Lynch, Spanish American Revolutions, 347–351. Rodríguez, Independence of Spanish America, 245. Kinsbruner, Independence in Spanish America, 142–143. - Lynch, Spanish American Revolutions, 347–349. - The Argentine President promotes Juana Azurduy to General in the Argentine Army.www.szmm.gov.hu/download.php?ctag=download&docID=14380 - ” O’Connor, Mothers Making Latin America”, 26-27. - Lynch, Spanish American Revolutions, 342–343. Kinsbruner, Independence in Spanish America, 146–152. - Lynch, Spanish American Revolutions, 351–352. Kinsbruner, Independence in Spanish America, 145–146, 152–153. - Rodríguez, Independence of Spanish America, 3–5, 213, 239. Kinsbruner states, "[I]n Mexico between 1820 and 1835 a larger percentage of adult males were permitted to vote than was the case in the United States, Great Britain, or France." Independence in Spanish America, 90. - Lynch, Spanish American Revolutions, 341–342, 352–355. Rodríguez, Independence of Spanish America, 219–222, 240–244. Kinsbruner, Independence in Spanish America, 143–144. - Andrews, George Reid. "Spanish American independence: A structural analysis." Latin American Perspectives (1985): 105-132. online - Andrien, Kenneth J. and Lyman L. Johnson. The Political Economy of Spanish America in the Age of Revolution, 1750–1850. Albuquerque, University of New Mexico Press, 1994. ISBN 978-0-8263-1489-5 - Anna, Timothy.. Spain & the Loss of Empire. Lincoln, University of Nebraska Press, 1983. ISBN 978-0-8032-1014-1 - Archer, Christon I., ed.. The Wars of Independence in Spanish America. Willmington, SR Books, 2000. ISBN 0-8420-2469-7 - Brading, D.A. The First America: The Spanish Monarchy, Creole Patriots and the Liberal State, 1492–1867. Cambridge University Press, 1991. ISBN 0-521-44796-8 - Chasteen, John Charles. Americanos: Latin America's Struggle for Independence. Oxford University Press, 2008. ISBN 978-0-19-517881-4 - Costeloe, Michael P. Response to Revolution: Imperial Spain and the Spanish American Revolutions, 1810–1840. Cambridge University Press, 1986. ISBN 978-0-521-32083-2 - Domínguez, Jorge I. Insurrection or Loyalty: The Breakdown of the Spanish American Empire. Cambridge, Harvard University Press, 1980. ISBN 978-0-674-45635-8 - Graham, Richard. Independence in Latin America: A Comparative Approach (2nd edition). McGraw-Hill, 1994. ISBN 0-07-024008-6 - Harvey, Robert. "Liberators: Latin America`s Struggle For Independence, 1810–1830". John Murray, London (2000). ISBN 0-7195-5566-3 - Higgins, James (editor). The Emancipation of Peru: British Eyewitness Accounts, 2014. Online at https://sites.google.com/site/jhemanperu - Humphreys, R. A., and John Lynch (editors). The Origins of the Latin American Revolutions, 1808–1826. New York, Alfred A. Knopf, 1965. - Kinsbruner, Jay. The Spanish-American Independence Movement. (Krieger Publishing Company, 1976). ISBN 978-0-88275-428-4 - Kinsbruner, Jay. Independence in Spanish America: Civil Wars, Revolutions, and Underdevelopment (2nd ed. University of New Mexico Press, 2000). ISBN 0-8263-2177-1 - Lynch, John. Caudillos in Spanish America, 1800–1850. Oxford, Clarendon Press, 1992. ISBN 0-19-821135-X - Lynch, John. The Spanish American Revolutions, 1808–1826 (2nd edition). New York, W. W. Norton & Company, 1986. ISBN 0-393-95537-0 - Lynch, John, ed. Latin American Revolutions, 1808-1826: Old and New World Origins (1995) 424pp; essays by scholrs - Rodríguez O., Jaime E. The Independence of Spanish America. Cambridge University Press, 1998. ISBN 0-521-62673-0 - Brown, Matthew. Adventuring through Spanish Colonies: Simón Bolívar, Foreign Mercenaries and the Birth of New Nations. Liverpool University Press, 2006. ISBN 1-84631-044-X - Hasbrouck, Alfred. Foreign Legionaries in the Liberation of Spanish South America. New York: Octagon Books, 1969. - Kaufman, William W. British Policy and the Independence of Latin America, 1804–1828. New Haven, Yale University Press, 1951. - Robertson, William Spence. France and Latin American Independence. (1939) - Rodríguez, Moises Enrique. Freedom's Mercenaries: British Volunteers in the Wars of Independence of Latin America, 2 vols. Lanham, Hamilton Books, University Press of America, 2006. ISBN 978-0-7618-3438-0 - Whitaker, Arthur P. The United States and the Independence of Latin America, 1800–1830. Baltimore, Johns Hopkins University Press, 1941. - Hensel, Silke. "Was There an Age of Revolution in Latin America?: New Literature on Latin American Independence." Latin American Research Review (2003) 38#3 pp: 237-249. online - Uribe, Victor M. "The Enigma of Latin American Independence: Analyses of the Last Ten Years," Latin American Research Review (1997) 32#1 pp. 236-255 in JSTOR
Let’s learn what is Integration before understanding the concept of Integration by Substitution. The integration of a function f(x) is given by F(x) and it is represented by: ∫f(x)dx = F(x) + C Here R.H.S. of the equation means integral of f(x) with respect to x. - F(x) is called anti-derivative or primitive. - f(x) is called the integrand. - dx is called the integrating agent. - C is called constant of integration or arbitrary constant. - x is the variable of integration. The anti-derivatives of basic functions are known to us. The integrals of these functions can be obtained readily. But this integration technique is limited to basic functions and in order to determine the integrals of various functions, different methods of integration are used. Among these methods of integration let us discuss integration by substitution. Integration by Substitution Method In this method of integration by substitution, any given integral is transformed into a simple form of integral by substituting the independent variable by others. Take for example an equation having an independent variable in x, i.e. ∫sin (x3).3x2.dx———————–(i), In the equation given above the independent variable can be transformed into another variable say t. Substituting x3 = t ———————-(ii) Differentiation of above equation will give- 3x2.dx = dt ———————-(iii) Substituting the value of (ii) and (iii) in (i), we have ∫sin (x3).3x2.dx = ∫sin t . dt Thus the integration of the above equation will give ∫sin t . dt= -cos t + c Again putting back the value of t from equation (ii), we get ∫sin (x3).3x2.dx = -cos x3 + c The General Form of integration by substitution is: ∫ f(g(x)).g'(x).dx = f(t).dt, where t = g(x) Usually the method of integration by substitution is extremely useful when we make a substitution for a function whose derivative is also present in the integrand. Doing so, the function simplifies and then the basic formulas of integration can be used to integrate the function. When to Use Integration by Substitution Method? In calculus, the integration by substitution method is also known as the “Reverse Chain Rule” or “U-Substitution Method”. We can use this method to find an integral value when it is set up in the special form. It means that the given integral is of the form: ∫ f(g(x)).g'(x).dx = f(u).du Here, first, integrate the function with respect to the substituted value (f(u)), and finish the process by substituting the original function g(x). Integration by Substitution Example To understand this concept better, let us look into the examples. Find the integration of Let t = tan-1x …… (1) dt = (1/ 1+x2 ) . dx I = ∫ et . dt = et + C …….(2) Substituting the value of (1) in (2), we have I = etan-1x + C. This is the required integration for the given function. Integrate 2x cos (x2 – 5) with respect to x . I = ∫2xcos(x2 – 5).dx Let x2 – 5 = t …..(1) 2x.dx = dt Substituting these values, we have I = ∫cos(t).dt = sin t + c …..(2) Substituting the value of 1 in 2, we have = sin (x2 – 5) + C This is the required integration for the given function. To learn more about integration by substitution please download Byju’s- The Learning App.
To understand gravitational waves must glimpse the structure of space-time defined by Albert Einstein in his theory of general relativity (1916). Einstein has linked three dimensions of space and one dimension of time, in a same tissue of space-time. This four-dimensional tissue "looks" on a trampoline area distended by the mass of planets, stars and galaxies. This distortion, compression or curvature of space-time in three dimensions is that we feel like gravity. In other words, a planet like the Earth is in orbit simply because it follows the curves of hummocky spatial tissue and distorted by the presence of the Sun and other planets of solar system. Thus gravitational waves (OG) are deformations of the structure of space-time that propagate as waves at the speed of light. They reflect the dynamics of space-time under the effect of rapid movements of ordinary matter while electromagnetic waves (photons) are produced by the movement of electric charges. In summary, gravitational waves or waves of curvature are produced by the acceleration of masses. Only the most relativistic objects in the universe, those who are extremely massive as black holes or neutron stars, can "shake" slightly space-time if they are accelerated. For example, two black holes 2 or neutron stars of the order of a few solar masses which rotate around each other, generate gravitational waves. But contrary to electromagnetic waves, gravitational waves interact very little with matter, they travel the cosmos without being absorbed, making them invisible in the electromagnetic pictures of our telescopes. nota: electromagnetic waves (radio, IR, optical, UV, X and gamma) are perturbations in the electromagnetic field, which propagate in space-time, while the gravitational waves are waves of space-time itself. In addition to detecting the variation (distance between peaks and troughs) is extremely low, the frequency to be detected, is a very low frequency. If we want to measure 10 000 km (radius of the wave), the variation of the gravitational wave generated by a black hole in our galaxy, the detectors must be able to observe a change in the wavelength of the size of an atom, that is 10-10 meter. Moreover, these variations are extremely rare in our Galaxy, we must substantially improve our current detectors and get those variations in other galaxies. Although not detected, scientists know they exist in the Universe. In 1975, radio astronomers Russell Hulse and Joseph Hooton Alan Taylor (Nobel Prize in Physics 1993) discovered, PSR B1913 +16 in the constellation Aquila, a binary pulsar (two neutron stars) with exceptional orbital characteristics. Indeed the two star orbit around one another in 7.75 hours in an extremely small volume of the order of 1.1 (periastron) to 4.8 times (apoastre) the radius of the sun. The small acceleration of the orbital period of this massive system (2 objects rotate faster and faster) and the shortening of the radius of the orbit (loss of 3 mm per orbit) has demonstrated the existence of gravitational waves. According to the theory of general relativity the orbit of a binary system is slowly modified by the emission of gravitational waves. Over thirty years Taylor and his colleagues made measurements which correspond exactly to the theory. For several other binary pulsars have confirmed the results of Taylor. The measures do not capable of detecting the energy of gravitational waves but are indirect evidence of the effects of gravitational waves emitted by a system. Video: representation of gravitational waves or bending waves generated by two black holes or two pulsars (neutron stars). A pulsar is a neutron star, extremely dense, density of an atomic nucleus, hence its name. Its giant magnetic field rotates about the axis of rotation to the rotation frequency of the star, some rotate in a millisecond, magnetic beam ejects particles which generate radio waves. They are cosmic beacons. When two black holes or two neutron stars rotate around each other, objects distort space-time and this deformation causes small gravitational waves like on video. Detection supposed of a gravitational wave A South Pole Telescope called BICEP-2 (Background Imaging of Cosmic Extragalactic Polarization 2) has enabled scientists to analyze the polarization of the light emitted by the early universe. This detection confirms Einstein's general relativity predicted the existence of gravitational waves, as a shiver of space-time caused by a large displacement of masses. How gravitational waves can be detected when they do not interact with matter? The telescope has detected a subtle property of the cosmic microwave background discovered in 1964, the famous primitive radiation Big Bang old 13,800,000,000 years. BICEP measured the polarization to large-scale microwave radiation. Only primordial gravitational waves can print such a model only if they were amplified by inflation. What is inflation? The distribution of matter in space is too uniform to be due solely Big bang. In the 1970s, cosmologists have imagined a sudden expansion of the universe, they called inflation. This inflation took place from the first second after Big Bang. Only inflation can amplify sufficiently primordial gravitational wave signal to make it detectable. Scientists with BICEP-2 specifically were looking to measure the polarization of the cosmic microwave background, i.e. the orientation of the electric field in the sky. They were looking for a specific type of polarization called B-modes, a model of vortex in the direction to polarized light from the ancient universe. In theory, this swirling pattern of polarization (see picture opposite) can only be created from gravitational waves. This is what BICEP-2 was found. " This is a very clean signature of these gravity waves," said Stanford physicist Kent Irwin. "But because of the importance of these results, they should be viewed with skepticism ... there is, in this condition, oddities in the results that may be disturbing ... I'm looking forward to seeing these results confirmed or reversed by other experiences, " said David Spergel, a professor of astrophysics at Princeton University. Indeed, the measurement is so difficult to do, it could easily be contaminated. Collaboration with other space telescopes such as Planck, should publish the results soon on the polarization of the cosmic microwave background. Further experiments are working on similar goals, which may support or go against BICEP-2. The June 5, 2014, at the Congress of the American Astronomical Society, David Spergel announced that the polarization B-Mode detected by BICEP2 could instead be the result of light scattering in the dust between the stars of the Milky Way. If the primordial gravitational waves are so anxiously hoped, is that they can confirm that inflation has occurred. Image: Gravitational waves of cosmic inflation, interpreted in the radiation of the cosmic microwave background image collected by a telescope experience BICEP-2 (Background Imaging of Cosmic Extragalactic Polarization) at the South Pole. Scientists estimate that the polarization or orientation wave, light is proof in the form of a signature called B-mode polarization or swirling pattern of polarization. This wave is represented on the image by small black lines guided by the whirlwind. The color indicates small temperature fluctuations in the cosmic microwave background that correspond to density fluctuations in the early universe.
Area is the space covered by a two-dimensional shape. The two dimensions are often expressed in terms of length and height or length and width. Multiplying length by height gives the area of rectangular objects. If these objects are measured in feet, the area is in square feet.Know More Two perpendicular measurements are needed to calculate the area of most shapes. On a rectangle, measuring the sides can give you the dimensions you need because one side can be the length and the adjacent side can be the height. On many other shapes, measuring one or more flat sides and then the height (or greatest distance taken with a perpendicular line) often works. For rectangles, multiply the length (or base) by the width (or height). The units are also multiplied. For example, if you measured a wall at 12 feet long by 8 feet high, the area would be 96 square feet. Use the appropriate formula for the shape. For example, the area of a triangle is the height multiplied by the base and divided by two; the area of a circle is the radius squared multiplied by pi. If the shape is irregular, you may have to calculate the area in segments. If you were working with units other than feet, convert the area to square feet. For example, if you used inches to measure the sides of the object, divide the resulting measurement in square inches by 144 to obtain the area in square feet. The surface area of a cube is the total sum of the area of the six sides that cover it, and it is calculated by using the formula A = 6a^2. Measure the area of one side of the cube. Multiply that number by 6 to arrive at the total surface area of the cube.Full Answer > To calculate the area of a circle, square the radius and multiply by pi (approximately 3.14). This enables you to find the area of any circle.Full Answer > To find the area of a triangle without a right angle, you multiply one-half the base by the height. Obtuse triangles are included in this group.Full Answer > Surface area is calculated by adding up all of the areas of the shapes that make up the surface of the object. The formula for calculating surface area differs with the type of object. It is vital to ensure that all units are the same for measuring the areas.Full Answer >
“A picture is worth a thousand words” is an English language adage meaning that complex and sometimes multiple ideas can be conveyed by a single still image, which conveys its meaning or essence more effectively than a mere verbal description (Wikipedia, 2020a, link). The origin of this statement can be traced back to the early 1900s. “In conveying information, an appropriate combination of pictures and words can far exceed either alone.” (David Moursund; American computer education leader; 1936-). The combination of mathematics and computers provides tools for storing, representing, analyzing, and communicating certain types of information. This and the next two newsletters focus on the use of graphical and other representations of historical data and information. Computers and mathematics enter this discussion in four major ways. The Web and the Internet together are an interactive, multimedia, storage, processing, distribution, and retrieval system that is steadily growing in the use of artificial intelligence. I strongly believe all teachers have a responsibility to help their students learn to make effective use of resources on the Web and Internet that are pertinent to the disciplines they are teaching. Students need to learn how to access, understand, interpret, and use historical data and information in text and/or represented graphically. They also need to learn to create their own graphic representations that will make effective use of this data and information. Computers are very useful aids in the latter task. I have long been amused by the quotation “A picture is worth a thousand words” that I used to start the newsletter. An amusing aside. Nowadays, most photographs are produced by use of digital cameras. A digital camera creates a computer representation of the scene being photographed. We can compare the amount of computer memory needed to store a picture with the number of words (of text) that can be stored in the same amount of computer memory. The digital cameras built into today’s widely sold smartphone vary somewhat in terms of the number of pixels (picture elements, colored dots) in a picture. However, roughly speaking, such a camera uses six million bytes (48 million bits) of storage space (or more) for one picture (Wikipedia, 2020b, link). We can use the estimate that an average English word is about five characters long, and that one character can be stored in one byte of computer memory, or five bytes per average word. This means that one such digital picture uses enough computer memory to store about 1.2 million words. Thus, the statement that a picture is worth a thousand words is archaic. In terms of computer storage requirements, an online picture is worth well over a million words! Moreover, it is easy to edit a picture stored in computer memory. This can be done at the level of a single pixel, just by changing binary bits representing the color of that pixel. This situation is one important aspect of the fake news problem. Incidentally, notice that the content given above mentioned both bits and bytes. These are standard words used in talking about computer memory. A bit is the amount of memory needed to store a binary bit—a 0 or a 1. A byte is the amount of memory needed to store 8 bits. Thus, a byte can store any one of the 256 different binary numbers 00000000, 00000001, 00000010, …, 11111110, 11111111. A computer also uses these binary codes to represent the letters and punctuation marks in an alphabet-based written language. When you are writing using a word processor, each key you press sends an eight-bit code (one byte of data) to a computer memory location. When I use the term picture, you probably think of a photograph or painting. In the next section, a pie chart in considered to be a type of picture. One commonly used example of visual representation is a pie chart or pie graph, with different segments representing parts of a whole. For most people, a pie chart is easier to understand than a table of numbers. The basic idea of parts of a whole applies to all pie charts, no matter what the topic of the information being presented. The example below is based on the idea of percentage, with the whole (collection of data) being 100%. Figure 1. An example of a pie chart (Math Is Fun, n.d., link). In Figure 1, the circle (the pie) is divided into 100 equal-sized slices, each being 1% of the total area of the pie. The necessary computations, as well as drawing, coloring, and labeling the pie chart, were done by a computer. A third-grade student can easily learn to use a computer to create such pie charts. The student can see the differences in sizes of the segments and begin to gain insight into percentages. However, using by hand tools is another matter. It is common to measure angles in degrees, where a degree in 1/360 of a circle. A protractor can be used for this purpose. The arithmetic involved is beyond the typical third grader. Moreover, it is common to draw circles using a compass that has a very sharp point along with a pencil or pen. A compass is potentially a dangerous weapon! A number of websites offer free software for creating pie charts. The National Center for Educational Statistics provides both the software and a tutorial (NCES, n.d., link). If you are not familiar with such software, I recommend you go to the site and follow the tutorial to make a pie chart. Notice how easy it is to experiment with colors and with the size of the circle. Some of this may be strictly trial and error as you are creating a pie chart that is pleasing to your eye and fits well into the document you are writing. In the United States, it is common for fourth grade students to study the history of their state. Suppose, for a moment, that you are a fourth-grade teacher who is teaching the history of the state where your students are living. Of what possible use would it be for you and your students to know how to read and create pie charts? Well, perhaps a bit of history will provide us with some help in answering this question. Quoting from Who Made That Pie Chart (Greenbaum & Rubinstein, 4/20/2012, link): William Playfair — a businessman, engineer and economics writer from Scotland — created the first known pie chart in 1801. Seeking to illustrate the Turkish Empire’s landholdings for his statistical breviary on the European nation-states, Playfair sliced a circle into three wedges whose sizes were determined by land area. According to a paper called “No Humble Pie: The Origins and Usage of a Statistical Chart” by Ian Spence, professor of psychology at the University of Toronto, Asia’s wedge was the largest and colored green, to indicate that it was a maritime power. Europe was made red to indicate it was a “land power.” It’s not clear why Playfair made Africa yellow, or what inspired him to make a circle graph at all. But, Spence writes, it was “the first pie chart to display empirical proportions and to differentiate the component parts by color.” I find that very interesting. The pie chart was invented only 200 years ago! Compare that with paper-and-pencil arithmetic invented more than 5,000 years ago, and the 2,300-year-old content of the geometry course you had in secondary school. And, what about algebra? Quoting from The Father of Algebra (Aljazeera, 10/20/2015, link): Hmm. When you took your first algebra course, did you learn this bit of history? Each discipline of study has a history.' In summary, when your students are learning about pie charts, they are leaning math that is much newer (more recently developed) that most of the rest of the math they will encounter in their education. When you introduce your students to the use of computers to create pie charts, you are helping them to learn about computer graphics, a very important and useful component of the discipline of Computer and Information Science. Graphical ways of representing information are taught in the school mathematics curriculum. However, students can learn to read graphs without knowing the mathematics needed to create them. At a more advanced mathematics level, similar statements hold for the mathematics of probability and statistics, a topic that will be covered in a later newsletter in this series. You know that reading and writing are two different—but closely related—areas of study. Think about this in terms of reading (getting information from) pie charts and writing (creating) pie charts). A similar type of analysis applies to graphs. Mathematicians and others have developed many different types of graphs to represent data. This section introduces line graphs, bar graphs, area graphs, and X-Y plot graphs as examples of representing historical data. Figure 2 shows a bar graph produced using free software available on the Web (Intellspot, n.d., link). In terms of representing information, it is much like a pie chart. Both a pie chart and a bar graph make use of color and size to distinguish among the elements being presented. However, a pie chart represents parts of a whole while a bar graph does not. In contrast, a vertically-oriented or a horizontally-oriented bar chart is probably better than a pie chart for representing a large number of pieces of information. Also, I find it much easier to compare the sizes of the bars (the vertical heights of the bars) than to compare the sizes of the wedges in a pie chart. Figure 2. An example of a bar graph. Figure 3 shows two ways—a table and a line graph— to represent a collection of data (Math Goodies, n.d., link). Both the table and the line graph represent the same historical data. Sarah bought a new car in 2001 for $24,000. Cars depreciate in value over the years. The table and the line graph both show this depreciation. Most people can more easily detect the somewhat linear decrease in value by viewing the line graph than they can by viewing the table. However, the table can easily be expanded to contain a column showing the depreciation for each year, as in Figure 4 below. I find this example puzzling. The decrease in value in the first year appears to be much too small. Perhaps Sarah only owed the car for part of 2001? I am also suspicious of the large decrease in value in the years 2004, 2005, and 2006. I know that older cars depreciate much more slowly than newer cars. My conclusion is that this is made up data (fake data) —not a set of true data. Warning, warning, warning! Just because you find data and information on the Web does not mean that it is correct. All teachers who have their students retrieve information from the Web have the responsibility of teaching their students that each has a personal responsibility to check for correctness of the information being retrieved. As I have just illustrated, often this can be done by using one’s accumulated knowledge and common sense. Figure 5 is a line graph from the National Center for Educational Statistics (NCES, n.d., link). This graph summarizes a considerable collection of historical data and information. Notice that both the horizontal and the vertical scales are equal interval scales, one being years and the other being thousands of dollars. Color is used to distinguish between inflated (blue) and uninflated (orange, constant) dollars. The blue and orange lines indicate increasing expenditures, but also indicate that quite a bit of this increase has been due to inflation. Thus, $2 thousand per pupil in 1977-78 dollars corresponds to $6.1 thousand per pupil in 2002-2003 dollars. This same data and information could have been represented in a table. But, the graphical representation makes it much easier to “see” and understand what has actually been going on. Figure 5. Line graph of two sets of data. Figure 6 shows an area graph (Stacked Area Chart, 6/1/19, link). Area graphs (area charts) are very similar to line graphs. They can be used to track changes over time for one or more groups. Area graphs may be useful when you are tracking the changes in two or more related groups. Figure 6. Area graph of five states’ gross product over time. X-Y Plot Graphs Figure 7 provides an example of an X-Y Plot Graph (Matange, 10/4/2016, link). X-Y plot graphs are also called scatter graphs. They are used to determine relationships between two different things. The x-axis is used to measure one event (or variable) and the y-axis is used to measure the other. Statistically, if both variables increase at the same time, they have a positive relationship. If one variable decreases while the other increases, they have a negative relationship. Sometimes the variables do not follow any pattern and have no relationship. This topic will be covered in a later newsletter that discusses correlation and causality. Figure 7. X-Y Plot graph of weight versus height. There are many more types of graphs. An article by Samantha Lile lists 44 types (Lile, n.d., link). She emphasizes the importance of selecting a type of graph that will best fit the data you want to communicate as well as best meet the needs of your audience. Computers are a major change agent in our schools. In each discipline of study, computers are affecting the instructional content, pedagogy, and assessment. The discipline of history provides an interesting case study. The world’s accumulated collection of readily accessible historical knowledge is growing rapidly. This includes substantial knowledge about individual people that currently is being gathered through the social networking systems, through online and offline stores, and through the Web search engines. This latter collection of data and information is not often available to the general public, and much of it is being used for commercial and other purposes. Interactive multimedia instructional systems are growing in capability through the use of Artificial Intelligence, and these systems are changing pedagogy. I frequently watch Public Broadcasting System programs and other educational programs on my TV. I often think of their instructional value versus the type of history education I received when I was in K-12 schools. Wow! What a powerful aid to the content and pedagogy aspects of history education! History teachers also face the issue of authentic assessment. What do we want students to memorize as they are being prepared for an adult life where they will wear a wrist device that provides voice input to and output from the Web? In addition, the issue of open book versus closed book exams has now become an issue of open Web-connected computers versus closed computer exams. The next two newsletters will continue to explore these and related issues. Aljazeera (10/20/2015). Al-Khwarizmi: The father of algebra. Retrieved 2/8/2020 from https://www.aljazeera.com/programmes/science-in-a-golden-age/2015/10/al-khwarizmi-father-algebra-151019144853758.html. Greenbaum, H., & and Rubinstein, D. (4/20/2012). Who made that pie chart? The New York Times Magazine. Retrieved 2/8/2020 from https://www.nytimes.com/2012/04/22/magazine/who-made-that-pie-chart.html. Intellspot (n.d.). Best free statistical graphing software. Retrieved 2/8/2020 from http://intellspot.com/best-free-graphing-software/. Lile, S. (n.d.). Forty-four types of graphs perfect for every top industry. visme. Retrieved 2/9/2020 from https://visme.co/blog/types-of-graphs/. Matange, S. (10/4/2016). Getting started with SGPLOT - Part 1 - Scatter plot. SAS. Retrieved 2/11/2020 from https://blogs.sas.com/content/graphicallyspeaking/2016/10/04/getting-started-sgplot-part-1-scatter-plot/, Math Goodies (n.d.). Line graphs. Lessons. Retrieved 2/9/2020 from https://www.mathgoodies.com/lessons/graphs/line. Math Is Fun (n.d.). Pie chart. Retrieved 2/9/2020 from https://www.mathsisfun.com/data/pie-charts.html Moursund, D. (2018a). The fourth R (Second edition). Eugene, OR: Information Age Education. Retrieved 2/8/2020 from http://iae-pedia.org/The_Fourth_R_(Second_Edition). Download the Microsoft Word file from http://i-a-e.org/downloads/free-ebooks-by-dave-moursund/307-the-fourth-r-second-edition.html. Download the PDF file from http://i-a-e.org/downloads/free-ebooks-by-dave-moursund/308-the-fourth-r-second-edition-1.html. See the Spanish edition, La cuarta R, below. Moursund, D. (2018b). La cuarta R. Eugene, OR: Information Age Education. Retrieved 2/8/2020 from http://iae-pedia.org/La_Cuarta_R_(Segunda_Edici%C3%B3n). NCES (n.d.) Kids’ zone. National Center for Education Statistics. Retrieved 2/8/2020 from https://nces.ed.gov/nceskids/graphing/Classic/pie.asp. Stacked Area Chart (n.d.). Great Lakes gross state product. Retrieved 2/11/2020 from https://docs.devexpress.com/AspNet/15910/aspnet-webforms-controls/chart-control/concepts/creating-charts/2d-chart-types/stacked-area-chart. Wikipedia (2020a). A picture is worth a thousand words. Retrieved 2/3/2020 from https://en.wikipedia.org/wiki/A_picture_is_worth_a_thousand_words. Wikipedia (2020b). 8-bit color. Retrieved 2/6/2020 from https://en.wikipedia.org/wiki/8-bit_color. David Moursund is an Emeritus Professor of Education at the University of Oregon, and editor of the IAE Newsletter. His professional career includes founding the International Society for Technology in Education (ISTE) in 1979, serving as ISTE’s executive officer for 19 years, and establishing ISTE’s flagship publication, Learning and Leading with Technology (now published by ISTE as Empowered Learner). He was the major professor or co-major professor for 82 doctoral students. He has presented hundreds of professional talks and workshops. He has authored or coauthored more than 60 academic books and hundreds of articles. Many of these books are available free online. See http://iaepedia.org/David_Moursund_Books . In 2007, Moursund founded Information Age Education (IAE). IAE provides free online educational materials via its IAE-pedia, IAE Newsletter, IAE Blog, and IAE books. See http://iaepedia.org/Main_Page#IAE_in_a_Nutshell . Information Age Education is now fully integrated into the 501(c)(3) non-profit corporation, Advancement of Globally Appropriate Technology and Education (AGATE) that was established in 2016. David Moursund is the Chief Executive Officer of AGATE. Email: firstname.lastname@example.orgReader Comments
Subject specialists have created NCERT Solutions for Class 11 Maths Chapter 15 Statistics, which includes thorough solutions for reference. All of the questions from the textbook’s exercises are answered here. Students can use these answers to help them prepare for their exams. The NCERT Solutions for Class 11 provide useful solutions for improving conceptual knowledge. The solutions are carefully solved using student-friendly terms while still adhering to the norms that must be followed when solving NCERT Solutions for Class 11. Practicing these answers can be incredibly advantageous not only in terms of exams but also in terms of helping Class 11 pupils perform well in upcoming competitive exams. The approaches for answering have been given special consideration to stay on target while not deviating from the intended answer. Because time is so important in exams, excellent time management when answering questions is essential for getting the best results. NCERT Solutions for Class 11 Maths Chapter 15 Statistics \|Exercise 15.1 Solutions\| \|Exercise 15.2 Solutions\| \|Exercise 15.3 Solutions\| \|Miscellaneous Exercise Solutions\| NCERT Solutions for Class 11 Maths Chapter 15 – Statistics Following are the topics and subtopics of Class 11 Maths Chapter 15 – Statistics of NCERT Solution: With examples, this section discusses the concepts of central tendency, mean, and median [during even and odd numbers of observations]. It introduces the notion of dispersion measurement. Measures of central tendency are values that cluster around the middle or center of the distribution. The three terms are mean, median, and mode. The mean can be used to find the average marks earned by students in a class. The median can be used to obtain the midway value when calculating student height. 15.2 Measures of Dispersion This section defines dispersion metrics, including the range, quartile deviation, mean deviation, and standard deviation. The relationship between measures of central tendency and measures of dispersion is explained by measures of dispersion. The spread of data, for example, indicates how well the mean represents the data. The mean is not reflective of data if the spread is big. The range, its formula, and an example are all defined in this section. Using the set’s highest and minimum values, the range calculates the variability of scores. Batsman A’s range in a cricket match is 121 – 0 = 121, while Batsman B’s range is 52 – 22 = 30. Range A is greater than Range B. As a result, in the instance of A, the scores are dispersed, whereas in the case of B, they are near together. 15.4 Mean Deviation The formula for mean deviation is defined in this section. Biologists can utilize the notion of mean deviation to compare different animal weights and determine what is a healthy weight. 15.4.1 Mean deviation for ungrouped data In this section, we’ll go through how to calculate the mean deviation for ungrouped data. Find the mean, standard deviations, and absolute deviations, then use the mean deviation formula to calculate the solution. 15.4.2 Mean deviation for grouped data This section explains how to calculate mean deviation for discrete and continuous frequency distributions using solved examples. 15.4.3 Limitations of mean deviation - If there is a lot of variability in a series, the median will not be a good representation of the data. As a result, the mean deviation derived around such a median is unreliable. - The mean deviation from the mean is not very specific if the sum of deviations from the mean is bigger than the sum of deviations from the median. - Further algebraic handling of the obtained absolute mean deviation is not possible. It can’t be utilized as a reliable indicator of dispersion. 15.5 Variance and standard deviation 15.5.1 Standard Deviation The definitions of variance and standard deviation, as well as formulas and solved examples for discrete and continuous frequency distributions, are covered in this section. A group of pupils completed a scientific test. On calculation, the test had an average score of 85 percent. The teacher looked up the standard deviation of the other results and discovered that it was relatively minimal, implying that the majority of the kids scored very close to 85 percent. 15.5.2 Shortcut method to find Variance and standard deviation With a few illustrations, this section explains how to calculate the standard deviation in a simplified approach. 15.6 Analysis of Frequency Distributions This section explains how to compare the variability of two series with the same mean, coefficient of variation, and a small number of solved issues. Following are the concepts discussed in Chapter 15 Statistics: - The difference between the maximum and least value of the given data is defined as the range. - If two series have the same mean, the one with the lower standard deviation is more stable and less scattered. - The variance is unaffected by adding or subtracting a positive value from each data point in the data collection. The solutions provide alternative tactics and clarifications for dealing with problems, allowing the learner to feel confident when taking the first term exam. In addition, dealing with a variety of confusing difficulties improves pupils’ knowledge of mathematical abilities. The answers cover all of the important questions that a student must know in order to appear for the term – I test. The Infinity Learn topic specialists who wrote the NCERT Solutions for Class 11 Maths based on the most recent CBSE Syllabus 2021-22 have a thorough understanding of the question paper setting and the distribution of marks throughout the chapters. Frequently Asked Questions on NCERT Solutions for Class 11 Maths Chapter 15 - Explain the term standard deviation which is discussed in Chapter 15 of the NCERT Solutions for Class 11 Maths? - What are the topics discussed in Chapter 15 of NCERT Solutions for Class 11 Maths? - Are the NCERT Solutions for Class 11 Maths Chapter 15 helpful for the students? 1. Explain the term standard deviation which is discussed in Chapter 15 of the NCERT Solutions for Class 11 Maths? The measurement of variation or deviation of a given set of values is called standard deviation. The range is determined by the standard deviation level. Students are urged to obtain the PDF of solutions accessible at Infinity study to better grasp this term. The solutions are presented in a chapter-by-chapter and exercise-by-exercise format to meet the needs of pupils. 2. What are the topics discussed in Chapter 15 of NCERT Solutions for Class 11 Maths? Following are the topics covered in Chapter 15 of NCERT Solutions for Class 11 Maths: - Methods of Dispersion - Mean Deviation - Variance and Standard Deviation - Analysis of Frequency Distributions 3. Are the NCERT Solutions for Class 11 Maths Chapter 15 helpful for the students? The NCERT Solutions for Class 11 Maths Chapter 15 provide precise explanations in easy language to assist students in performing well in the first term exams. According to the new CBSE Syllabus 2021-22, the step-by-step manner of solving issues gives students a clear concept of the marks’ weightage. Students will be able to identify their areas of weakness and try to improve them in order to improve their academic performance.
African Americans are an ethnic group of Americans with total or partial ancestry from any of the black racial groups of Africa. The term refers to descendants of enslaved black people who are from the United States. Black and African Americans constitute the third largest racial and ethnic group in the United States. Most African Americans are descendants of enslaved peoples within the boundaries of the present United States. On average, African Americans are of West/Central African and European descent, some have Native American ancestry. According to U. S. Census Bureau data, African immigrants do not self-identify as African American; the overwhelming majority of African immigrants identify instead with their own respective ethnicities. Immigrants from some Caribbean, Central American and South American nations and their descendants may or may not self-identify with the term. African-American history starts in the 16th century, with peoples from West Africa forcibly taken as slaves to Spanish America, in the 17th century with West African slaves taken to English colonies in North America. After the founding of the United States, black people continued to be enslaved, the last four million black slaves were only liberated after the Civil War in 1865. Due to notions of white supremacy, they were treated as second-class citizens; the Naturalization Act of 1790 limited U. S. citizenship to whites only, only white men of property could vote. These circumstances were changed by Reconstruction, development of the black community, participation in the great military conflicts of the United States, the elimination of racial segregation, the civil rights movement which sought political and social freedom. In 2008, Barack Obama became the first African American to be elected President of the United States; the first African slaves arrived via Santo Domingo to the San Miguel de Gualdape colony, founded by Spanish explorer Lucas Vázquez de Ayllón in 1526. The marriage between Luisa de Abrego, a free black domestic servant from Seville and Miguel Rodríguez, a white Segovian conquistador in 1565 in St. Augustine, is the first known and recorded Christian marriage anywhere in what is now the continental United States. The ill-fated colony was immediately disrupted by a fight over leadership, during which the slaves revolted and fled the colony to seek refuge among local Native Americans. De Ayllón and many of the colonists died shortly afterwards of an epidemic and the colony was abandoned; the settlers and the slaves who had not escaped returned to Haiti, whence. The first recorded Africans in British North America were "20 and odd negroes" who came to Jamestown, Virginia via Cape Comfort in August 1619 as indentured servants; as English settlers died from harsh conditions and more Africans were brought to work as laborers. An indentured servant would work for several years without wages; the status of indentured servants in early Virginia and Maryland was similar to slavery. Servants could be bought, sold, or leased and they could be physically beaten for disobedience or running away. Unlike slaves, they were freed after their term of service expired or was bought out, their children did not inherit their status, on their release from contract they received "a year's provision of corn, double apparel, tools necessary", a small cash payment called "freedom dues". Africans could raise crops and cattle to purchase their freedom. They raised families, married other Africans and sometimes intermarried with Native Americans or English settlers. By the 1640s and 1650s, several African families owned farms around Jamestown and some became wealthy by colonial standards and purchased indentured servants of their own. In 1640, the Virginia General Court recorded the earliest documentation of lifetime slavery when they sentenced John Punch, a Negro, to lifetime servitude under his master Hugh Gwyn for running away. In the Spanish Florida some Spanish married or had unions with Pensacola, Creek or African women, both slave and free, their descendants created a mixed-race population of mestizos and mulattos; the Spanish encouraged slaves from the southern British colonies to come to Florida as a refuge, promising freedom in exchange for conversion to Catholicism. King Charles II of Spain issued a royal proclamation freeing all slaves who fled to Spanish Florida and accepted conversion and baptism. Most went to the area around St. Augustine, but escaped slaves reached Pensacola. St. Augustine had mustered an all-black militia unit defending Spain as early as 1683. One of the Dutch African arrivals, Anthony Johnson, would own one of the first black "slaves", John Casor, resulting from the court ruling of a civil case; the popular conception of a race-based slave system did not develop until the 18th century. The Dutch West India Company introduced slavery in 1625 with the importation of eleven black slaves into New Amsterdam. All the colony's slaves, were freed upon its surrender to the British. Massachusetts was the first British colony to recognize slavery in 1641. In 1662, Virginia passed a law that children of enslaved women took the status of the mother, rather than that of the father, as under English common law; this principle was called partus sequitur ventrum. By an act of 1699, the colony ordered all free blacks deported defining as slaves all people of African descent who remained in the c Virtual International Authority File The Virtual International Authority File is an international authority file. It is a joint project of several national libraries and operated by the Online Computer Library Center. Discussion about having a common international authority started in the late 1990s. After a series of failed attempts to come up with a unique common authority file, the new idea was to link existing national authorities; this would present all the benefits of a common file without requiring a large investment of time and expense in the process. The project was initiated by the US Library of Congress, the German National Library and the OCLC on August 6, 2003; the Bibliothèque nationale de France joined the project on October 5, 2007. The project transitioned to being a service of the OCLC on April 4, 2012; the aim is to link the national authority files to a single virtual authority file. In this file, identical records from the different data sets are linked together. A VIAF record receives a standard data number, contains the primary "see" and "see also" records from the original records, refers to the original authority records. The data are available for research and data exchange and sharing. Reciprocal updating uses the Open Archives Initiative Protocol for Metadata Harvesting protocol; the file numbers are being added to Wikipedia biographical articles and are incorporated into Wikidata. VIAF's clustering algorithm is run every month; as more data are added from participating libraries, clusters of authority records may coalesce or split, leading to some fluctuation in the VIAF identifier of certain authority records. Authority control Faceted Application of Subject Terminology Integrated Authority File International Standard Authority Data Number International Standard Name Identifier Wikipedia's authority control template for articles Official website VIAF at OCLC New York City The City of New York called either New York City or New York, is the most populous city in the United States. With an estimated 2017 population of 8,622,698 distributed over a land area of about 302.6 square miles, New York is the most densely populated major city in the United States. Located at the southern tip of the state of New York, the city is the center of the New York metropolitan area, the largest metropolitan area in the world by urban landmass and one of the world's most populous megacities, with an estimated 20,320,876 people in its 2017 Metropolitan Statistical Area and 23,876,155 residents in its Combined Statistical Area. A global power city, New York City has been described as the cultural and media capital of the world, exerts a significant impact upon commerce, research, education, tourism, art and sports; the city's fast pace has inspired the term New York minute. Home to the headquarters of the United Nations, New York is an important center for international diplomacy. Situated on one of the world's largest natural harbors, New York City consists of five boroughs, each of, a separate county of the State of New York. The five boroughs – Brooklyn, Manhattan, The Bronx, Staten Island – were consolidated into a single city in 1898; the city and its metropolitan area constitute the premier gateway for legal immigration to the United States. As many as 800 languages are spoken in New York, making it the most linguistically diverse city in the world. New York City is home to more than 3.2 million residents born outside the United States, the largest foreign-born population of any city in the world. In 2017, the New York metropolitan area produced a gross metropolitan product of US$1.73 trillion. If greater New York City were a sovereign state, it would have the 12th highest GDP in the world. New York is home to the highest number of billionaires of any city in the world. New York City traces its origins to a trading post founded by colonists from the Dutch Republic in 1624 on Lower Manhattan. The city and its surroundings came under English control in 1664 and were renamed New York after King Charles II of England granted the lands to his brother, the Duke of York. New York served as the capital of the United States from 1785 until 1790, it has been the country's largest city since 1790. The Statue of Liberty greeted millions of immigrants as they came to the U. S. by ship in the late 19th and early 20th centuries and is an international symbol of the U. S. and its ideals of liberty and peace. In the 21st century, New York has emerged as a global node of creativity and entrepreneurship, social tolerance, environmental sustainability, as a symbol of freedom and cultural diversity. Many districts and landmarks in New York City are well known, with the city having three of the world's ten most visited tourist attractions in 2013 and receiving a record 62.8 million tourists in 2017. Several sources have ranked New York the most photographed city in the world. Times Square, iconic as the world's "heart" and its "Crossroads", is the brightly illuminated hub of the Broadway Theater District, one of the world's busiest pedestrian intersections, a major center of the world's entertainment industry. The names of many of the city's landmarks and parks are known around the world. Manhattan's real estate market is among the most expensive in the world. New York is home to the largest ethnic Chinese population outside of Asia, with multiple signature Chinatowns developing across the city. Providing continuous 24/7 service, the New York City Subway is the largest single-operator rapid transit system worldwide, with 472 rail stations. Over 120 colleges and universities are located in New York City, including Columbia University, New York University, Rockefeller University, which have been ranked among the top universities in the world. Anchored by Wall Street in the Financial District of Lower Manhattan, New York has been called both the most economically powerful city and the leading financial center of the world, the city is home to the world's two largest stock exchanges by total market capitalization, the New York Stock Exchange and NASDAQ. In 1664, the city was named in honor of the Duke of York. James's older brother, King Charles II, had appointed the Duke proprietor of the former territory of New Netherland, including the city of New Amsterdam, which England had seized from the Dutch. During the Wisconsinan glaciation, 75,000 to 11,000 years ago, the New York City region was situated at the edge of a large ice sheet over 1,000 feet in depth; the erosive forward movement of the ice contributed to the separation of what is now Long Island and Staten Island. That action left bedrock at a shallow depth, providing a solid foundation for most of Manhattan's skyscrapers. In the precolonial era, the area of present-day New York City was inhabited by Algonquian Native Americans, including the Lenape, whose homeland, known as Lenapehoking, included Staten Island; the first documented visit into New York Harbor by a European was in 1524 by Giovanni da Verrazzano, a Florentine explorer in the service of the French crown. He named it Nouvelle Angoulême. A Spanish expedition led by captain Estêvão Gomes, a Portuguese sailing for Emperor Charles V, arrived in New York Harbor in January 1525 and charted the mouth of the Hudson River, which he named Río de San Antonio. The Padrón Rea Système universitaire de documentation The système universitaire de documentation or SUDOC is a system used by the libraries of French universities and higher education establishments to identify and manage the documents in their possession. The catalog, which contains more than 10 million references, allows students and researcher to search for bibliographical and location information in over 3,400 documentation centers, it is maintained by the Bibliographic Agency for Higher Education. Official website Washington, D. C. formally the District of Columbia and referred to as Washington or D. C. is the capital of the United States. Founded after the American Revolution as the seat of government of the newly independent country, Washington was named after George Washington, first President of the United States and Founding Father; as the seat of the United States federal government and several international organizations, Washington is an important world political capital. The city is one of the most visited cities in the world, with more than 20 million tourists annually; the signing of the Residence Act on July 16, 1790, approved the creation of a capital district located along the Potomac River on the country's East Coast. The U. S. Constitution provided for a federal district under the exclusive jurisdiction of the U. S. Congress, the District is therefore not a part of any state; the states of Maryland and Virginia each donated land to form the federal district, which included the pre-existing settlements of Georgetown and Alexandria. The City of Washington was founded in 1791 to serve as the new national capital. In 1846, Congress returned the land ceded by Virginia. Washington had an estimated population of 702,455 as of July 2018, making it the 20th most populous city in the United States. Commuters from the surrounding Maryland and Virginia suburbs raise the city's daytime population to more than one million during the workweek. Washington's metropolitan area, the country's sixth largest, had a 2017 estimated population of 6.2 million residents. All three branches of the U. S. federal government are centered in the District: Congress and the U. S. Supreme Court. Washington is home to many national monuments, museums situated on or around the National Mall; the city hosts 177 foreign embassies as well as the headquarters of many international organizations, trade unions, non-profit, lobbying groups, professional associations, including the World Bank Group, the International Monetary Fund, the Organization of American States, AARP, the National Geographic Society, the Human Rights Campaign, the International Finance Corporation, the American Red Cross. A locally elected mayor and a 13‑member council have governed the District since 1973. However, Congress may overturn local laws. D. C. residents elect a non-voting, at-large congressional delegate to the House of Representatives, but the District has no representation in the Senate. The District receives three electoral votes in presidential elections as permitted by the Twenty-third Amendment to the United States Constitution, ratified in 1961. Various tribes of the Algonquian-speaking Piscataway people inhabited the lands around the Potomac River when Europeans first visited the area in the early 17th century. One group known as the Nacotchtank maintained settlements around the Anacostia River within the present-day District of Columbia. Conflicts with European colonists and neighboring tribes forced the relocation of the Piscataway people, some of whom established a new settlement in 1699 near Point of Rocks, Maryland. In his Federalist No. 43, published January 23, 1788, James Madison argued that the new federal government would need authority over a national capital to provide for its own maintenance and safety. Five years earlier, a band of unpaid soldiers besieged Congress while its members were meeting in Philadelphia. Known as the Pennsylvania Mutiny of 1783, the event emphasized the need for the national government not to rely on any state for its own security. Article One, Section Eight, of the Constitution permits the establishment of a "District as may, by cession of particular states, the acceptance of Congress, become the seat of the government of the United States". However, the Constitution does not specify a location for the capital. In what is now known as the Compromise of 1790, Alexander Hamilton, Thomas Jefferson came to an agreement that the federal government would pay each state's remaining Revolutionary War debts in exchange for establishing the new national capital in the southern United States. On July 9, 1790, Congress passed the Residence Act, which approved the creation of a national capital on the Potomac River; the exact location was to be selected by President George Washington, who signed the bill into law on July 16. Formed from land donated by the states of Maryland and Virginia, the initial shape of the federal district was a square measuring 10 miles on each side, totaling 100 square miles. Two pre-existing settlements were included in the territory: the port of Georgetown, founded in 1751, the city of Alexandria, founded in 1749. During 1791–92, Andrew Ellicott and several assistants, including a free African American astronomer named Benjamin Banneker, surveyed the borders of the federal district and placed boundary stones at every mile point. Many of the stones are still standing. A new federal city was constructed on the north bank of the Potomac, to the east of Georgetown. On September 9, 1791, the three commissioners overseeing the capital's construction named the city in honor of President Washington; the federal district was named Columbia, a poetic name for the United States in use at that time. Congress held its first session in Washington on November 17, 1800. Congress passed the District of Columbia Organic Act of 1801 that organized the District and placed the entire territory under the exclusive control of the federal Howard University is a private, federally chartered black university in Washington, D. C, it is categorized by the Carnegie Foundation as a research university with higher research activity and is accredited by the Middle States Commission on Higher Education. From its outset Howard has been open to people of all sexes and races. Howard offers more than 120 areas leading to undergraduate and professional degrees. Howard is classified as a Tier 1 national university and ranks second among HBCUs by U. S. News & World Report. Howard is the only HBCU ranked in the top 40 on the Bloomberg Businessweek college rankings; the Princeton Review ranked the school of business first in opportunities for minority students and in the top five for most competitive students. The National Law Journal ranked the law school among the top 25 in the nation for placing graduates at the most successful law firms. Howard has produced four Rhodes Scholars between 1986 and 2017. Between 1998 and 2018, Howard University produced two Marshall Scholars, eleven Truman Scholars, seventy Fulbright Scholars, a Schwarzman Scholar and twenty-two Pickering Fellows. Howard produces the most black doctorate recipients of any university. Shortly after the end of the American Civil War, members of The First Congregational Society of Washington considered establishing a theological seminary for the education of African-American clergymen. Within a few weeks, the project expanded to include a provision for establishing a university. Within two years, the University consisted of the Colleges of Liberal Medicine; the new institution was named for General Oliver Otis Howard, a Civil War hero, both the founder of the University and, at the time, Commissioner of the Freedmen's Bureau. Howard served as President of the University from 1869–74. U. S. Congress chartered Howard on March 2, 1867, much of its early funding came from endowment, private benefaction, tuition. (In the 20th and 21st centuries an annual congressional appropriation, administered by the U. S. Department of Education, funds Howard University and Howard University Hospital After five years of being an institution, Howard University became the place of education for over 150,000 freed slaves. Many improvements were made on campus. Howard Hall was made a dormitory for women. From 1926-1960, Howard University's first African-American presideant, Dr. Mordecai Wyatt Johnson, Sr. reigned. The Great Depression years of the 1930s brought hardship to campus. Despite appeals from Eleanor Roosevelt, Howard saw its budget cut below Hoover administration levels during the Presidency of Franklin D. Roosevelt. Howard University has played an important role in American history and the Civil Rights Movement on a number of occasions. Alain Locke, Chair of the Department of Philosophy and first African American Rhodes Scholar, authored The New Negro, which helped to usher in the Harlem Renaissance. Ralph Bunche, the first Nobel Peace Prize winner of African descent, served as chair of the Department of Political Science. Beginning in 1942, Howard University students pioneered the "stool-sitting" technique of occupying stools at a local cafeteria which denied service to African Americans blocking other customers waiting for service. This tactic was to play a prominent role in the Civil Rights Movement. By January 1943, students had begun to organize regular sit-ins and pickets at cigar stores and cafeterias around Washington, D. C. which refused to serve them because of their race. These protests continued until the fall of 1944. Stokely Carmichael known as Kwame Toure, a student in the Department of Philosophy and the Howard University School of Divinity, coined the term "Black Power" and worked in Lowndes County, Alabama as a voting rights activist. Historian Rayford Logan served as chair of the Department of History. E. Franklin Frazier served as chair of the Department of Sociology. Sterling Allen Brown served as chair of the Department of English; the first sitting president to speak at Howard was Calvin Coolidge in 1924. His graduation speech was entitled, "The Progress of a People," and highlighted the accomplishments to date of the blacks in America since the Civil War, his concluding thought was, "We can not go out from this place and occasion without refreshment of faith and renewal of confidence that in every exigency our Negro fellow citizens will render the best and fullest measure of service whereof they are capable." In 1965, President Lyndon B. Johnson delivered a speech to the graduating class at Howard, where he outlined his plans for civil rights legislation and endorsed aggressive affirmative action to combat the effects of years of segregation of blacks from the nation's economic opportunities. At the time, the Voting Rights bill was still pending in the House of Representatives. In 1975 the historic Freedman's Hospital closed after 112 years of use as Howard University College of Medicine's primary teaching hospital. Howard University Hospital opened that same year and continues to be used as Howard University College of Medicine's primary teaching hospital with service to the surrounding community. In 1989, Howard gained national attention when students rose up in protest against the appointment of then-Republican National Committee Chairman Lee Atwater as a new member of the university's board of trustees. Student activists disrupted Howard's 122nd anniversary celebrations, occupied the university's Administration building. Within days, both Atwater and Howard's President, James E. Cheek, resigned. In April 2007, the head of the faculty senate called for the ouster of Howard University President H. Patrick Swygert, saying the school was in a state of crisis and it was time to end "an intolerable condition of incompetence Yale University is a private Ivy League research university in New Haven, Connecticut. Founded in 1701, it is the third-oldest institution of higher education in the United States and one of the nine Colonial Colleges chartered before the American Revolution. Chartered by Connecticut Colony, the "Collegiate School" was established by clergy to educate Congregational ministers, it moved to New Haven in 1716 and shortly after was renamed Yale College in recognition of a gift from British East India Company governor Elihu Yale. Restricted to theology and sacred languages, the curriculum began to incorporate humanities and sciences by the time of the American Revolution. In the 19th century, the college expanded into graduate and professional instruction, awarding the first Ph. D. in the United States in 1861 and organizing as a university in 1887. Its faculty and student populations grew after 1890 with rapid expansion of the physical campus and scientific research. Yale is organized into fourteen constituent schools: the original undergraduate college, the Yale Graduate School of Arts and Sciences and twelve professional schools. While the university is governed by the Yale Corporation, each school's faculty oversees its curriculum and degree programs. In addition to a central campus in downtown New Haven, the university owns athletic facilities in western New Haven, a campus in West Haven and forest and nature preserves throughout New England; the university's assets include an endowment valued at $29.4 billion as of October 2018, the second largest endowment of any educational institution in the world. The Yale University Library, serving all constituent schools, holds more than 15 million volumes and is the third-largest academic library in the United States. Yale College undergraduates follow a liberal arts curriculum with departmental majors and are organized into a social system of residential colleges. All members of the Faculty of Arts and Sciences—and some members of other faculties—teach undergraduate courses, more than 2,000 of which are offered annually. Students compete intercollegiately as the Yale Bulldogs in the NCAA Division I – Ivy League. As of October 2018, 61 Nobel laureates, 5 Fields Medalists and 3 Turing award winners have been affiliated with Yale University. In addition, Yale has graduated many notable alumni, including five U. S. Presidents, 19 U. S. Supreme Court Justices, 31 living billionaires and many heads of state. Hundreds of members of Congress and many U. S. diplomats, 78 MacArthur Fellows, 247 Rhodes Scholars and 119 Marshall Scholars have been affiliated with the university. Its wealth and influence have led to Yale being reported as amoungst the most prestigious universities in the United States. Yale traces its beginnings to "An Act for Liberty to Erect a Collegiate School", passed by the General Court of the Colony of Connecticut on October 9, 1701, while meeting in New Haven; the Act was an effort to create an institution to train ministers and lay leadership for Connecticut. Soon thereafter, a group of ten Congregational ministers, Samuel Andrew, Thomas Buckingham, Israel Chauncy, Samuel Mather, Rev. James Noyes II, James Pierpont, Abraham Pierson, Noadiah Russell, Joseph Webb, Timothy Woodbridge, all alumni of Harvard, met in the study of Reverend Samuel Russell in Branford, Connecticut, to pool their books to form the school's library. The group, led by James Pierpont, is now known as "The Founders". Known as the "Collegiate School", the institution opened in the home of its first rector, Abraham Pierson, today considered the first president of Yale. Pierson lived in Killingworth; the school moved to Saybrook and Wethersfield. In 1716, it moved to Connecticut. Meanwhile, there was a rift forming at Harvard between its sixth president, Increase Mather, the rest of the Harvard clergy, whom Mather viewed as liberal, ecclesiastically lax, overly broad in Church polity; the feud caused the Mathers to champion the success of the Collegiate School in the hope that it would maintain the Puritan religious orthodoxy in a way that Harvard had not. In 1718, at the behest of either Rector Samuel Andrew or the colony's Governor Gurdon Saltonstall, Cotton Mather contacted the successful Boston born businessman Elihu Yale to ask him for financial help in constructing a new building for the college. Through the persuasion of Jeremiah Dummer, Elihu "Eli" Yale, who had made a fortune through trade while living in Madras as a representative of the East India Company, donated nine bales of goods, which were sold for more than £560, a substantial sum at the time. Cotton Mather suggested that the school change its name to "Yale College".. Meanwhile, a Harvard graduate working in England convinced some 180 prominent intellectuals that they should donate books to Yale; the 1714 shipment of 500 books represented the best of modern English literature, science and theology. It had a profound effect on intellectuals at Yale. Undergraduate Jonathan Edwards discovered John Locke's works and developed his original theology known as the "new divinity". In 1722 the Rector and six of his friends, who had a study group to discuss the new ideas, announced that they had given up Calvinism, become Arminians and joined the Church of England, they were returned to the colonies as missionaries for the Anglican faith. Thomas Clapp became president in 1745 and struggled to return the college to Calvinist orthodoxy, but he did not close the library. Other students found Deist books in the library. Yale was swept up by the great intellectual movements of the peri
What is algorithm and flowchart with example? Algorithm and flowchart are the powerful tools for learning programming. An algorithm is a step-by-step analysis of the process, while a flowchart explains the steps of a program in a graphical way. Algorithm and flowcharts helps to clarify all the steps for solving the problem. How should a beginner learn a flowchart? How to plan and draw a basic flowchart - Define your purpose and scope. - Identify the tasks in chronological order. - Organize them by type and corresponding shape, such as process, decision, data, inputs or outputs. - Draw your chart, either sketching by hand or using a program such as Lucidchart. What is an example of a simple algorithm? Common examples include: the recipe for baking a cake, the method we use to solve a long division problem, the process of doing laundry, and the functionality of a search engine are all examples of an algorithm. How do I write an algorithm? There are many ways to write an algorithm….An Algorithm Development Process - Step 1: Obtain a description of the problem. This step is much more difficult than it appears. - Step 2: Analyze the problem. - Step 3: Develop a high-level algorithm. - Step 4: Refine the algorithm by adding more detail. - Step 5: Review the algorithm. How do you write an algorithm for beginners? How do I create an algorithm flowchart in Word? How to make a Flowchart in Word - Open a blank document in Word. - Add shapes. To begin adding shapes to your flowchart in Word, you have two options. - Add text. Add text to a SmartArt graphic by clicking the filler text and begin typing. - Add lines. - Format shapes and lines. What is C++ flowchart? Flowchart is a diagrammatic representation of sequence of logical steps of a program. Flowcharts use simple geometric shapes to depict processes and arrows to show relationships and process/data flow. What are the 4 algorithms? Algorithm types we will consider include: - Simple recursive algorithms. - Backtracking algorithms. - Divide and conquer algorithms. - Dynamic programming algorithms. - Greedy algorithms. - Branch and bound algorithms. - Brute force algorithms. - Randomized algorithms. What are the 3 algorithms? Types of Algorithm - Recursive Algorithm. This is one of the most interesting Algorithms as it calls itself with a smaller value as inputs which it gets after solving for the current inputs. - Divide and Conquer Algorithm. - Dynamic Programming Algorithm. - Greedy Algorithm. - Brute Force Algorithm. - Backtracking Algorithm. What is the difference between a flowchart and an algorithm? The definition of algorithm and flowchart. What is Algorithm? How can algorithms be used to describe a flowchart? – Document flowcharts, showing controls over a document-flow through a system – Data flowcharts, showing controls over a data-flow in a system – System flowcharts, showing controls at a physical or resource level – Program flowchart, showing the controls in a program within a system What are the steps of algorithm? – In OTTs, algorithms comprise actions and calculations to automate a search or offer recommendations. – This can impact the art of cinema which was Academy Award-winning filmmaker Martin Scorcese’s main concern. – Some OTTs predicted this and have opted for a more human approach. Which is difficult to design flowchart or algorithem? There are no stringent rules are implemented in the algorithms while the flowchart is abode by predefined rules. Errors and bugs are easily detected in the algorithm as compared to the flow charts. Flow charts are simple to create.
Hey , I still do not understand how to determine the activation energy for the reverse reaction when you are given the activation energy for the forward reaction and the delta h for the reaction. This state is also known as an activated complex. The energy values points on the hyper-surface along the reaction coordinate result in a 1-D energy surface a line and when plotted against the reaction coordinate energy vs reaction coordinate gives what is called a reaction coordinate diagram or energy profile. For any reaction to proceed, the starting material must have enough energy to cross over an energy barrier. In this equation, k is the rate constant for the reaction, Z is a proportionality constant that varies from one reaction to another, E a is the activation energy for the reaction, R is the ideal gas constant in joules per mole kelvin, and T is the temperature in kelvin. The reaction coordinate is described by its parameters, which are frequently given as a composite of several geometric parameters, and can change direction as the reaction progresses so long as the smallest energy barrier or activation energy Ea is traversed. The saddle point represents the highest energy point lying on the reaction coordinate connecting the reactant and product; this is known as the transition state. The methods for describing the potential energy are broken down into a classical mechanics interpretation and a interpretation. Next focus on finding the energy level of the reactants or products and simply draw a double arrowhead line that connects the two energy levels, i. One guideline for drawing diagrams for complex reactions is the which says that a favored reaction proceeding from a reactant to an intermediate or from one intermediate to another or product is one which has the least change in nuclear position or electronic configuration. However, when more than one such barrier is to be crossed, it becomes important to recognize the highest barrier which will determine the rate of the reaction. Figure 5:Potential Energy Surface and Corresponding 2-D Reaction Coordinate Diagram derived from the plane passing through the minimum energy pathway between A and C and passing through B A reaction involving more than one elementary step has one or more intermediates being formed which, in turn, means there is more than one energy barrier to overcome. This is called kinetic control and the ratio of the products formed depends on the relative energy barriers leading to the products. If the starting material and product s are in equilibrium then their relative abundance is decided by the difference in free energy between them. These parameters are independent of each other. As the temperature of the system increases, the number of molecules that carry enough energy to react when they collide also increases. Figure 8: Reaction Coordinate Diagrams showing favorable or unfavorable and slow or fast reactions The relative stability of reactant and product does not define the feasibility of any reaction all by itself. The height of energy barrier is always measured relative to the energy of the reactant or starting material. Chemists use reaction coordinate diagrams as both an analytical and pedagogical aid for rationalizing and illustrating and events. So the activation energy for the reverse reaction is the sum of the enthalpy delta H and the activation energy Eact for the forward reaction. Not all reactions are reversible. However, if the two energy barriers for reactant-to-intermediate and intermediate-to-product transformation are nearly equal, then no complete equilibrium is established and steady state approximation is invoked to derive the kinetic rate expressions for such a reaction. Which statement describes the potential energy diagram of an endothermic reaction? Four criteria must be satisfied in order for something to be classified as catalyst. The vertical axis is the potential energy energy of the chemical bonds of the substances, and the horizontal axis is time. Depending on these parameters, a reaction can be favorable or unfavorable, fast or slow and reversible or irreversible, as shown in figure 8. In other words, a saddle point represents a transition state along the reaction coordinate. In principle, all elementary steps are reversible, but in many cases the equilibrium lies so much towards the product side that the starting material is effectively no longer observable or present in sufficient concentration to have an effect on reactivity. Reaction coordinate diagrams also give information about the equilibrium between a reactant or a product and an intermediate. The diagram represents the potential energy changes when a cold pack is activated. Thus, it can be said that the reactions involving dramatic changes in position of nuclei actually occur through a series of simple chemical reactions. I don't want the answer to become too long, so I won't go into too much detail about activation energy and threshold energy. In other words, the approximation allows the kinetic energy of the nuclei or movement of the nuclei to be neglected and therefore the nuclei repulsion is a constant value as static point charges and is only considered when calculating the total energy of the system. A reaction is in equilibrium when the rate of forward reaction is equal to the rate of reverse reaction. The third criterion is a consequence of the second; because catalysts are not consumed in the reaction, they can catalyze the reaction over and over again. The purpose of a catalyst is to alter the activation energy. For a forward reaction, the activation energy is equal to the difference between the threshold energy and the energy level of the reactants. Figure 13 shows a common way to illustrate the effect of an enzyme on a given biochemical reaction. The ground states are represented by local energy minima and the transition states by saddle points. So your activation energy is the distance between the potential energy of your reactants which is the flat line in the left of the diagram and the peak which is the highest point in the diagram. Explanation: 1 The potential energy diagram shows the evolution of the potential chemical energy of a process or reaction, measured as energy of the chemical bonds, including the reactants, the intermediates, the activated complex, and the products. When a reactant can form two different products depending on the reaction conditions, it becomes important to choose the right conditions to favor the desired product. A reaction can also be rendered irreversible if a subsequent, faster step takes place to consume the initial product s , or a gas is evolved in an open system. These changes in geometry of a molecule or interactions between molecules are dynamic processes which call for understanding all the forces operating within the system. If the forward reaction is endothermic, reactants will be lower in the energy diagram than products. This is important because the kinetic energy molecules carry when they collide is the principal source of the energy that must be invested in a reaction to get it started. The rate of a reaction depends on the temperature at which it is run. In subsequent steps, the activation energy is only from the intermediate to the next transition state. . Potential energy diagrams Chemical reactions involve a change in energy, usually a loss or gain of heat energy. Endothermic Reactions the reactants have less potential energy than do the products. In all three of these reactions the first step is the slow step because the activation energy from the reactants to the transition state is the highest. Modern Quantum Chemistry: Introduction to Advanced Electronic Structure Theory. Saddle point represents a maximum along only one direction that of the reaction coordinate and is a minimum along all other directions. An is a biological catalyst that increases the rate for many vital biochemical reactions. Positive catalysts increase the reaction rate and negative catalysts or inhibitors slow down a reaction and possibly cause the reaction not occur at all. I do know that the activation energy for the reverse reaction is larger than the forward reaction.
\|One of the Seven Basic Tools of Quality\| \|First described by\|\|Karl Pearson\| \|Purpose\|\|To roughly assess the probability distribution of a given variable by depicting the frequencies of observations occurring in certain ranges of values.\| A histogram is an accurate representation of the distribution of numerical data. It is an estimate of the probability distribution of a continuous variable (quantitative variable) and was first introduced by Karl Pearson. It differs from a bar graph, in the sense that a bar graph relates two variables, but a histogram relates only one. To construct a histogram, the first step is to "bin" the range of values—that is, divide the entire range of values into a series of intervals—and then count how many values fall into each interval. The bins are usually specified as consecutive, non-overlapping intervals of a variable. The bins (intervals) must be adjacent, and are often (but are not required to be) of equal size. If the bins are of equal size, a rectangle is erected over the bin with height proportional to the frequency—the number of cases in each bin. A histogram may also be normalized to display "relative" frequencies. It then shows the proportion of cases that fall into each of several categories, with the sum of the heights equaling 1. However, bins need not be of equal width; in that case, the erected rectangle is defined to have its area proportional to the frequency of cases in the bin. The vertical axis is then not the frequency but frequency density—the number of cases per unit of the variable on the horizontal axis. Examples of variable bin width are displayed on Census bureau data below. As the adjacent bins leave no gaps, the rectangles of a histogram touch each other to indicate that the original variable is continuous. Histograms give a rough sense of the density of the underlying distribution of the data, and often for density estimation: estimating the probability density function of the underlying variable. The total area of a histogram used for probability density is always normalized to 1. If the length of the intervals on the x-axis are all 1, then a histogram is identical to a relative frequency plot. A histogram can be thought of as a simplistic kernel density estimation, which uses a kernel to smooth frequencies over the bins. This yields a smoother probability density function, which will in general more accurately reflect distribution of the underlying variable. The density estimate could be plotted as an alternative to the histogram, and is usually drawn as a curve rather than a set of boxes. An alternative is the average shifted histogram, which is fast to compute and gives a smooth curve estimate of the density without using kernels. Histograms are sometimes confused with bar charts. A histogram is used for continuous data, where the bins represent ranges of data, while a bar chart is a plot of categorical variables. Some authors recommend that bar charts have gaps between the rectangles to clarify the distinction. The etymology of the word histogram is uncertain. Sometimes it is said to be derived from the Ancient Greek ἱστός (histos) – "anything set upright" (as the masts of a ship, the bar of a loom, or the vertical bars of a histogram); and γράμμα (gramma) – "drawing, record, writing". It is also said that Karl Pearson, who introduced the term in 1891, derived the name from "historical diagram". This is the data for the histogram to the right, using 500 items: \|−3.5 to -2.51\|\|9\| \|−2.5 to -1.51\|\|32\| \|−1.5 to -0.51\|\|109\| \|−0.5 to 0.49\|\|180\| \|0.5 to 1.49\|\|132\| \|1.5 to 2.49\|\|34\| \|2.5 to 3.49\|\|4\| The words used to describe the patterns in a histogram are: "symmetric", "skewed left" or "right", "unimodal", "bimodal" or "multimodal". It is a good idea to plot the data using several different bin widths to learn more about it. Here is an example on tips given in a restaurant. Here are a couple more examples: The U.S. Census Bureau found that there were 124 million people who work outside of their homes. Using their data on the time occupied by travel to work, the table below shows the absolute number of people who responded with travel times "at least 30 but less than 35 minutes" is higher than the numbers for the categories above and below it. This is likely due to people rounding their reported journey time. The problem of reporting values as somewhat arbitrarily rounded numbers is a common phenomenon when collecting data from people. This histogram shows the number of cases per unit interval as the height of each block, so that the area of each block is equal to the number of people in the survey who fall into its category. The area under the curve represents the total number of cases (124 million). This type of histogram shows absolute numbers, with Q in thousands. This histogram differs from the first only in the vertical scale. The area of each block is the fraction of the total that each category represents, and the total area of all the bars is equal to 1 (the fraction meaning "all"). The curve displayed is a simple density estimate. This version shows proportions, and is also known as a unit area histogram. In other words, a histogram represents a frequency distribution by means of rectangles whose widths represent class intervals and whose areas are proportional to the corresponding frequencies: the height of each is the average frequency density for the interval. The intervals are placed together in order to show that the data represented by the histogram, while exclusive, is also contiguous. (E.g., in a histogram it is possible to have two connecting intervals of 10.5–20.5 and 20.5–33.5, but not two connecting intervals of 10.5–20.5 and 22.5–32.5. Empty intervals are represented as empty and not skipped.) In a more general mathematical sense, a histogram is a function mi that counts the number of observations that fall into each of the disjoint categories (known as bins), whereas the graph of a histogram is merely one way to represent a histogram. Thus, if we let n be the total number of observations and k be the total number of bins, the histogram mi meets the following conditions: A cumulative histogram is a mapping that counts the cumulative number of observations in all of the bins up to the specified bin. That is, the cumulative histogram Mi of a histogram mj is defined as: There is no "best" number of bins, and different bin sizes can reveal different features of the data. Grouping data is at least as old as Graunt's work in the 17th century, but no systematic guidelines were given until Sturges's work in 1926. Using wider bins where the density of the underlying data points is low reduces noise due to sampling randomness; using narrower bins where the density is high (so the signal drowns the noise) gives greater precision to the density estimation. Thus varying the bin-width within a histogram can be beneficial. Nonetheless, equal-width bins are widely used. Some theoreticians have attempted to determine an optimal number of bins, but these methods generally make strong assumptions about the shape of the distribution. Depending on the actual data distribution and the goals of the analysis, different bin widths may be appropriate, so experimentation is usually needed to determine an appropriate width. There are, however, various useful guidelines and rules of thumb. The number of bins k can be assigned directly or can be calculated from a suggested bin width h as: The braces indicate the ceiling function. which takes the square root of the number of data points in the sample (used by Excel histograms and many others). Sturges' formula is derived from a binomial distribution and implicitly assumes an approximately normal distribution. It implicitly bases the bin sizes on the range of the data and can perform poorly if n < 30, because the number of bins will be small—less than seven—and unlikely to show trends in the data well. It may also perform poorly if the data are not normally distributed. The Rice Rule is presented as a simple alternative to Sturges's rule. Doane's formula is a modification of Sturges' formula which attempts to improve its performance with non-normal data. where is the estimated 3rd-moment-skewness of the distribution and where is the sample standard deviation. Scott's normal reference rule is optimal for random samples of normally distributed data, in the sense that it minimizes the integrated mean squared error of the density estimate. which is based on the interquartile range, denoted by IQR. It replaces 3.5σ of Scott's rule with 2 IQR, which is less sensitive than the standard deviation to outliers in data. Here, is the number of datapoints in the kth bin, and choosing the value of h that minimizes J will minimize integrated mean squared error. where and are mean and biased variance of a histogram with bin-width , and . A good reason why the number of bins should be proportional to is the following: suppose that the data are obtained as independent realizations of a bounded probability distribution with smooth density. Then the histogram remains equally "rugged" as tends to infinity. If is the "width" of the distribution (e. g., the standard deviation or the inter-quartile range), then the number of units in a bin (the frequency) is of order and the relative standard error is of order . Comparing to the next bin, the relative change of the frequency is of order provided that the derivative of the density is non-zero. These two are of the same order if is of order , so that is of order . This simple cubic root choice can also be applied to bins with non-constant width. \|Wikimedia Commons has media related to Histograms.\| \|Wikimedia Commons has media related to Histogram.\| \|Look up histogram in Wiktionary, the free dictionary.\| None of the audio/visual content is hosted on this site. All media is embedded from other sites such as GoogleVideo, Wikipedia, YouTube etc. Therefore, this site has no control over the copyright issues of the streaming media. All issues concerning copyright violations should be aimed at the sites hosting the material. 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With tuning instruments that can produce sustained tones, beats can be readily recognized. Tuning two tones to a unison will present a peculiar effect: when the two tones are close in pitch but not identical, the difference in frequency generates the beating. The volume varies like in a tremolo as the sounds alternately interfere constructively and destructively. As the two tones gradually approach unison, the beating slows down and may become so slow as to be imperceptible. As the two tones get further apart, their beat frequency starts to approach the range of human pitch perception, the beating starts to sound like a note, and a combination tone is produced. This combination tone can also be referred to as a missing fundamental, as the beat frequency of any two tones is equivalent to the frequency of their implied fundamental frequency. Mathematics and physics of beat tones This phenomenon is best known in acoustics or music, though it can be found in any linear system: "According to the law of superposition, two tones sounding simultaneously are superimposed in a very simple way: one adds their amplitudes". If a graph is drawn to show the function corresponding to the total sound of two strings, it can be seen that maxima and minima are no longer constant as when a pure note is played, but change over time: when the two waves are nearly 180 degrees out of phase the maxima of one wave cancel the minima of the other, whereas when they are nearly in phase their maxima sum up, raising the perceived volume. It can be proven (see List of trigonometric identities) that the envelope of the maxima and minima form a wave whose frequency is half the difference between the frequencies of the two original waves. Consider two sine waves of unit amplitude: If the two original frequencies are quite close (for example, a difference of approximately twelve hertz), the frequency of the cosine of the right side of the expression above, that is f1 − f2/, is often too low to be perceived as an audible tone or pitch. Instead, it is perceived as a periodic variation in the amplitude of the first term in the expression above. It can be said that the lower frequency cosine term is an envelope for the higher frequency one, i.e. that its amplitude is modulated. The frequency of the modulation is f1 + f2/, that is, the average of the two frequencies. It can be noted that every second burst in the modulation pattern is inverted. Each peak is replaced by a trough and vice versa. However, because the human ear is not sensitive to the phase of a sound, only its amplitude or intensity, only the magnitude of the envelope is heard. Therefore, subjectively, the frequency of the envelope seems to have twice the frequency of the modulating cosine, which means the audible beat frequency is: This can be seen on the adjacent diagram. A physical interpretation is that when the two waves are in phase and they interfere constructively. When it is zero, they are out of phase and interfere destructively. Beats occur also in more complex sounds, or in sounds of different volumes, though calculating them mathematically is not so easy. For a human ear to hear beat phenomena, the ratio of frequencies should be less than or else the brain perceives them as two different frequencies. Beating can also be heard between notes that are near to, but not exactly, a harmonic interval, due to some harmonic of the first note beating with a harmonic of the second note. For example, in the case of perfect fifth, the third harmonic (i.e. second overtone) of the bass note beats with the second harmonic (first overtone) of the other note. As well as with out-of tune notes, this can also happen with some correctly tuned equal temperament intervals, because of the differences between them and the corresponding just intonation intervals: see Harmonic series (music)#Harmonics and tuning. A binaural beat is an auditory illusion perceived when two different pure-tone sine waves, both with frequencies lower than 1500 Hz, with less than a 40 Hz difference between them, are presented to a listener dichotically (one through each ear). For example, if a 530 Hz pure tone is presented to a subject's right ear, while a 520 Hz pure tone is presented to the subject's left ear, the listener will perceive the auditory illusion of a third tone, in addition to the two pure-tones presented to each ear. The third sound is called a binaural beat, and in this example would have a perceived pitch correlating to a frequency of 10 Hz, that being the difference between the 530 Hz and 520 Hz pure tones presented to each ear. Binaural-beat perception originates in the inferior colliculus of the midbrain and the superior olivary complex of the brainstem, where auditory signals from each ear are integrated and precipitate electrical impulses along neural pathways through the reticular formation up the midbrain to the thalamus, auditory cortex, and other cortical regions. Some potential benefits of binaural beats therapy may include: reduced stress, reduced anxiety, increased focus, increased concentration, increased motivation, increased confidence, and deeper meditation. Musicians commonly use interference beats objectively to check tuning at the unison, perfect fifth, or other simple harmonic intervals. Piano and organ tuners even use a method involving counting beats, aiming at a particular number for a specific interval. The composer Alvin Lucier has written many pieces that feature interference beats as their main focus. Italian composer Giacinto Scelsi, whose style is grounded on microtonal oscillations of unisons, extensively explored the textural effects of interference beats, particularly in his late works such as the violin solos Xnoybis (1964) and L'âme ailée / L'âme ouverte (1973), which feature them prominently (note that Scelsi treated and notated each string of the instrument as a separate part, so that his violin solos are effectively quartets of one-strings, where different strings of the violin may be simultaneously playing the same note with microtonal shifts, so that the interference patterns are generated). Composer Phill Niblock's music is entirely based on beating caused by microtonal differences. - Levitin, Daniel (2006). This is Your Brain on Music: The Science of a Human Obsession. Penguin. p. 22. - Winckel, Fritz (1967). Music, Sound and Sensation: A Modern Exposition, p.134. Courier. ISBN 9780486165820. - "Interference beats and Tartini tones", Physclips, UNSW.edu.au. - "Acoustics FAQ", UNSW.edu.au. - Roberts, Gareth E. (2016). From Music to Mathematics: Exploring the Connections, p.112. JHU. 9781421419190. - Oster, G (October 1973). "Auditory beats in the brain". Scientific American. 229 (4): 94–102. Bibcode:1973SciAm.229d..94O. doi:10.1038/scientificamerican1073-94. PMID 4727697. - "What are binaural beats and how do they work?". Medical News Today. Retrieved September 2, 2019. - Campbell, Murray; Greated, Clive A.; and Myers, Arnold (2004). Musical Instruments: History, Technology, and Performance of Instruments of Western Music, p.26. Oxford. ISBN 9780198165040. "Listening for beats can be a useful method of tuning a unison, for example between two strings on a lute,..." - Java applet, MIT - Acoustics and Vibration Animations, D.A. Russell, Pennsylvania State University - A Java applet showing the formation of beats due to the interference of two waves of slightly different frequencies - Yet another interactive Java applet; also shows equation of combined waves, including phase angle. - Lissajous Curves: Interactive simulation of graphical representations of musical intervals, beats, interference, vibrating strings - Thaut, Michael H. (2005). Rhythm, music, and the brain : scientific foundations and clinical applications (1st in paperback ed.). New York, NY [u.a.]: Routledge. ISBN 978-0415973700. - Berger, Jonathan; Turow, Gabe, eds. (2011). Music, science, and the rhythmic brain : cultural and clinical implications. Routledge. ISBN 9780415890595.
OBJECTIVE 2 \| Describe the roles of categories, hierarchies, definitions, and prototypes in concept formation. OBJECTIVE 3 \| Compare algorithms and heuristics as problem-solving strategies, and explain how insight differs from both of them. OBJECTIVE 4 \| Contrast confirmation bias and fixation, and explain how they can interfere with effective problem solving. OBJECTIVE 5 \| Contrast the representative and availability heuristics, and explain how they can cause us to underestimate or ignore important information. OBJECTIVE 6 \| Describe the drawbacks and advantages of overconfidence in decision making. OBJECTIVE 7 \| Describe how others can use framing to elicit from us the answers they want. OBJECTIVE 8 \| Explain how our preexisting beliefs can distort our logic. OBJECTIVE 9 \| Describe the remedy for belief perseverance phenomenon. OBJECTIVE 10 \| Describe the smart thinker’s reaction to using intuition. OBJECTIVE 11 \| Describe the basic structural units of language. OBJECTIVE 12 \| Trace the course of language acquisition from the babbling stage through two-word stage. OBJECTIVE 13 \| Discuss Skinner’s and Chomsky’s contributions to the nature-nurture debate over how children acquire language, and explain why statistical learning and critical periods are important concepts in children’s language learning. OBJECTIVE 14 \| Summarize Whorf’s linguistic determinism hypothesis, and comment on its standing in contemporary psychology. OBJECTIVE 15 \| Discuss the value of thinking in images. OBJECTIVE 16 \| List five cognitive skills shared by the great apes and humans. OBJECTIVE 17 \| Outline the arguments for and against the idea that animals and humans share the capacity for language. Thinking and Language Chapter 10 Thinking and Language <ul><li>Thinking </li></ul><ul><ul><li>Concepts </li></ul></ul><ul><ul><li>Solving Problems </li></ul></ul><ul><ul><li>Making Decisions and Forming Judgments </li></ul></ul><ul><ul><li>Belief Bias </li></ul></ul> Thinking and Language <ul><li>Language </li></ul><ul><ul><li>Language Structure </li></ul></ul><ul><ul><li>Language Development </li></ul></ul><ul><li>Thinking & Language </li></ul><ul><ul><li>Language Influences Thinking </li></ul></ul><ul><ul><li>Thinking in Images </li></ul></ul> Thinking and Language <ul><li>Animal Thinking and Language </li></ul><ul><ul><li>Do Animals Think? </li></ul></ul><ul><ul><li>Do Animals Exhibit Language? </li></ul></ul><ul><ul><li>The Case of the Apes </li></ul></ul> Thinking <ul><li>Thinking, or cognition, refers to a process that involves knowing, understanding, remembering, and communicating. </li></ul> Cognitive Psychologists <ul><li>Thinking involves a number of mental activities, which are listed below. Cognitive psychologists study these in great detail. </li></ul><ul><li>Concepts </li></ul><ul><li>Problem solving </li></ul><ul><li>Decision making </li></ul><ul><li>Judgment formation </li></ul> Concept <ul><li>The mental grouping of similar objects, events, ideas, or people. There are a variety of chairs but their common features define the concept of a chair . </li></ul> Category Hierarchies <ul><li>We organize concepts into category hierarchies. </li></ul>Courtesy of Christine Brune Development of Concepts We form some concepts with definitions. For example, a triangle has three sides. Mostly, we form concepts with mental images or typical examples ( prototypes ). For example, a robin is a prototype of a bird, but a penguin is not. Triangle (definition) Bird (mental image) Daniel J. Cox/ Getty Images J. Messerschmidt/ The Picture Cube Categories Once we place an item in a category, our memory shifts toward the category prototype. A computer generated face that was 70 percent Caucasian led people to classify it as Caucasian. Courtesy of Oliver Corneille Problem Solving <ul><li>There are two ways to solve problems: </li></ul>Algorithms: Methodical, logical rules or procedures that guarantee solving a particular problem. Algorithms <ul><li>Algorithms, which are very time consuming, exhaust all possibilities before arriving at a solution. Computers use algorithms. </li></ul>S P L O Y O C H Y G If we were to unscramble these letters to form a word using an algorithmic approach, we would face 907,208 possibilities. Heuristics Heuristics are simple, thinking strategies that allow us to make judgments and solve problems efficiently. Heuristics are less time consuming, but more error-prone than algorithms. B2M Productions/Digital Version/Getty Images Heuristics <ul><li>Heuristics make it easier for us to use simple principles to arrive at solutions to problems. </li></ul>S P L O Y O C H Y G S P L O Y O C H G Y P S L O Y O C H G Y P S Y C H O L O G Y Put a Y at the end, and see if the word begins to make sense. Insight <ul><li>Insight involves a sudden novel realization of a solution to a problem. Humans and animals have insight. </li></ul>Grande using boxes to obtain food Insight <ul><li>Brain imaging and EEG studies suggest that when an insight strikes (the “Aha” experience), it activates the right temporal cortex (Jung-Beeman, 2004). The time between not knowing the solution and realizing it is 0.3 seconds. </li></ul>From Mark Jung-Beekman, Northwestern University and John Kounios, Drexel University Obstacles in Solving Problems <ul><li>Confirmation Bias: A tendency to search for information that confirms a personal bias. </li></ul>2 – 4 – 6 Rule: Any ascending series of numbers. 1 – 2 – 3 would comply. Ss had difficulty figuring out the rule due to a confirmation bias (Wason, 1960). Mental Set <ul><li>A tendency to approach a problem in a particular way, especially if that way was successful in the past. </li></ul> Functional Fixedness A tendency to think only of the familiar functions of an object. Problem: Tie the two ropes together. Use a screw driver, cotton balls and a matchbox. ? Functional Fixedness Use the screwdriver as a weight, and tie it to the end of one rope. Swing it toward the other rope to tie the knot. The inability to think of the screwdriver as a weight is functional fixedness. ? Using and Misusing Heuristics <ul><li>Two kinds of heuristics, representative heuristics and availability heuristics , have been identified by cognitive psychologists. </li></ul>Amos Tversky Daniel Kahneman Courtesy of Greymeyer Award, University of Louisville and the Tversky family Courtesy of Greymeyer Award, University of Louisville and Daniel Kahneman Representativeness Heuristic <ul><li>Judging the likelihood of things or objects in terms of how well they seem to represent, or match, a particular prototype. </li></ul>Probability that that person is a truck driver is far greater than an ivy league professor just because there are more truck drivers than such professors. If you meet a slim, short, man who wears glasses and likes poetry, what do you think his profession would be? An Ivy league professor or a truck driver? Availability Heuristic <ul><li>Why does our availability heuristic lead us astray? </li></ul><ul><li>Whatever increases the ease of retrieving information increases its perceived availability. </li></ul>How is retrieval facilitated? <ul><li>How recently we have heard about the event. </li></ul><ul><li>How distinct it is. </li></ul><ul><li>How correct it is. </li></ul> Making Decision & Forming Judgments <ul><li>Each day we make hundreds of judgments and decisions based on our intuition, seldom using systematic reasoning. </li></ul> Overconfidence <ul><li>Intuitive heuristics, confirmation of beliefs, and the inclination to explain failures increase our overconfidence. Overconfidence is a tendency to overestimate the accuracy of our beliefs and judgments. </li></ul>At a stock market, both the seller and the buyer may be confident about their decisions on a stock. Exaggerated Fear <ul><li>The opposite of having overconfidence is having an exaggerated fear about what may happen. Such fears may be unfounded. </li></ul><ul><li>The 9/11 attacks led to a decline in air travel due to fear. </li></ul>AP/ Wide World Photos Framing Decisions <ul><li>Decisions and judgments may be significantly affected depending upon how an issue is framed. </li></ul>Example: What is the best way to market ground beef — as 25% fat or 75% lean? Belief Bias <ul><li>The tendency of one’s preexisting beliefs to distort logical reasoning by making invalid conclusions. </li></ul>God is love. Love is blind Ray Charles is blind. Ray Charles is God. Anonymous graffiti Belief Perseverance <ul><li>Belief perseverance is the tendency to cling to our beliefs in the face of contrary evidence. </li></ul>If you see that a country is hostile, you are likely to interpret their ambiguous actions as a sign of hostility (Jervis, 1985). Perils & Powers of Intuition Intuition may be perilous if unchecked, but may also be extremely efficient and adaptive. Language <ul><li>Language, our spoken, written, or gestured work, is the way we communicate meaning to ourselves and others. </li></ul>Language transmits culture. M. & E. Bernheim/ Woodfin Camp & Associates Language Structure <ul><li>Phonemes: The smallest distinct sound unit in a spoken language. For example: </li></ul><ul><li>bat, has three phonemes b · a · t </li></ul><ul><li>chat, has three phonemes ch · a · t </li></ul> Language Structure <ul><li>Morpheme: The smallest unit that carries a meaning. It may be a word or part of a word. For example: </li></ul><ul><li>Milk = milk </li></ul><ul><li>Pumpkin = pump . kin </li></ul><ul><li>Unforgettable = un · for · get · table </li></ul> Structuring Language Phrase Sentence Meaningful units (290,500) … meat, pumpkin. Words Smallest meaningful units (100,000) … un, for . Morphemes Basic sounds (about 40) … ea, sh . Phonemes Composed of two or more words (326,000) … meat eater. Composed of many words (infinite) … She opened the jewelry box. Grammar <ul><li>Grammar is the system of rules in a language that enable us to communicate with and understand others. </li></ul>Grammar Syntax Semantics Semantics <ul><li>Semantics is the set of rules by which we derive meaning from morphemes, words, and sentences. For example: </li></ul>Semantic rule tells us that adding –ed to the word laugh means that it happened in the past. Syntax <ul><li>Syntax consists of the rules for combining words into grammatically sensible sentences. For example: </li></ul>In English, syntactical rule says that adjectives come before nouns; white house . In Spanish, it is reversed; casa blanca . Language Development <ul><li>Children learn their native languages much before learning to add 2+2. </li></ul><ul><li>We learn, on average (after age 1), 3,500 words a year, amassing 60,000 words by the time we graduate from high school. </li></ul>Time Life Pictures/ Getty Images When do we learn language? <ul><li>Babbling Stage: Beginning at 4 months, the infant spontaneously utters various sounds, like ah-goo . Babbling is not imitation of adult speech. </li></ul> When do we learn language? One-Word Stage: Beginning at or around his first birthday, a child starts to speak one word at a time and is able to make family members understand him. The word doggy may mean look at the dog out there . When do we learn language? <ul><li>Two-Word Stage: Before the 2nd year a child starts to speak in two-word sentences. This form of speech is called telegraphic speech because the child speaks like a telegram: “Go car,” means I would like to go for a ride in the car. </li></ul> When do we learn language? Longer phrases: After telegraphic speech, children begin uttering longer phrases ( Mommy get ball ) with syntactical sense, and by early elementary school they are employing humor. You never starve in the desert because of all the sand-which-is there. Explaining Language Development <ul><li>Operant Learning : Skinner (1957, 1985) believed that language development may be explained on the basis of learning principles such as association, imitation, and reinforcement. </li></ul> Explaining Language Development <ul><li>2. Inborn Universal Grammar: Chomsky (1959, 1987) opposed Skinner’s ideas and suggested that the rate of language acquisition is so fast that it cannot be explained through learning principles, and thus most of it is inborn. </li></ul> Explaining Language Development <ul><li>3. Statistical Learning and Critical Periods: Well before our first birthday, our brains are discerning word breaks by statistically analyzing which syllables in hap-py-ba-by go together. These statistical analyses are learned during critical periods of child development. </li></ul> Genes, Brain, & Language <ul><li>Genes design the mechanisms for a language, and experience modifies the brain. </li></ul>Michael Newman/ Photo Edit, Inc. Eye of Science/ Photo Researchers, Inc. David Hume Kennerly/ Getty Images Language & Age <ul><li>Learning new languages gets harder with age. </li></ul> Language & Thinking <ul><li>Language and thinking intricately intertwine. </li></ul>Rubber Ball/ Almay Language Influences Thinking <ul><li>Linguistic Determinism: Whorf (1956) suggested that language determines the way we think. For example, he noted that the Hopi people do not have the past tense for verbs. Therefore, the Hopi cannot think readily about the past. </li></ul> Language Influences Thinking <ul><li>When a language provides words for objects or events, we can think about these objects more clearly and remember them. It is easier to think about two colors with two different names (A) than colors with the same name (B) (Özgen, 2004). </li></ul> Word Power <ul><li>Increasing word power pays its dividends. It pays for speakers and deaf individuals who learn sign language. </li></ul> Linguistic Determinism Questioned <ul><li>Although people from Papua New Guinea do not use our words for colors and shapes, they still perceive them as we do (Rosch, 1974). </li></ul> Thinking in Images <ul><li>To a large extent thinking is language-based. When alone, we may talk to ourselves. However, we also think in images. </li></ul>2. When we are riding our bicycle. 1. When we open the hot water tap. We don’t think in words, when: Images and Brain <ul><li>Imagining a physical activity activates the same brain regions as when actually performing the activity. </li></ul>Jean Duffy Decety, September 2003 Language and Thinking <ul><li>Traffic runs both ways between language and thinking. </li></ul> <ul><li>Do animals have a language? </li></ul>Animals & Language Honey bees communicate by dancing. The dance moves clearly indicate the direction of the nectar. Do Animals Think? <ul><li>Common cognitive skills in humans and apes include the following: </li></ul><ul><li>Concept formation. </li></ul><ul><li>Insight </li></ul><ul><li>Problem Solving </li></ul><ul><li>Culture </li></ul><ul><li>Mind? </li></ul>African grey parrot assorts red blocks from green balls. William Munoz Insight <ul><li>Chimpanzees show insightful behavior when solving problems. </li></ul>Sultan uses sticks to get food. Problem Solving <ul><li>Apes are famous, much like us, for solving problems. </li></ul>Chimpanzee fishing for ants. Courtesy of Jennifer Byrne, c/o Richard Byrne, Department of Psychology, University of St. Andrews, Scotland Animal Culture <ul><li>Animals display customs and culture that are learned and transmitted over generations. </li></ul>Dolphins using sponges as forging tools. Chimpanzee mother using and teaching a young how to use a stone hammer. Copyright Amanda K Coakes Michael Nichols/ National Geographic Society Mental States <ul><li>Can animals infer mental states in themselves and others? </li></ul><ul><li>To some extent. Chimps and orangutans (and dolphins) used mirrors to inspect themselves when a researcher put paint spots on their faces or bodies. </li></ul> Do Animals Exhibit Language? <ul><li>There is no doubt that animals communicate. </li></ul><ul><li>Vervet monkeys, whales and even honey bees communicate with members of their species and other species. </li></ul>Rico (collie) has a 200-word vocabulary Copyright Baus/ Kreslowski The Case of Apes <ul><li>Chimps do not have a vocal apparatus for human-like speech (Hayes & Hayes,1951). Therefore, Gardner and Gardner (1969) used American Sign Language (ASL) to train Washoe, a chimp, who learned 182 signs by the age of 32. </li></ul> Gestured Communication <ul><li>Animals, like humans, exhibit communication through gestures. It is possible that vocal speech developed from gestures during the course of evolution. </li></ul> Sign Language <ul><li>American Sign Language (ASL) is instrumental in teaching chimpanzees a form of communication. </li></ul>When asked, this chimpanzee uses a sign to say it is a baby. Paul Fusco/ Magnum Photos Computer Assisted Language <ul><li>Others have shown that bonobo pygmy chimpanzees can develop even greater vocabularies and perhaps semantic nuances in learning a language (Savage-Rumbaugh, 1991). Kanzi and Panbanish developed vocabulary for hundreds of words and phrases. </li></ul>Copyright of Great Ape Trust of Iowa Criticism <ul><li>Apes acquire their limited vocabularies with a great deal of difficulty, unlike children who develop vocabularies at amazing rates. </li></ul><ul><li>Chimpanzees can make signs to receive a reward, just as a pigeon who pecks at the key receives a reward. However, pigeons have not learned a language. </li></ul><ul><li>Chimpanzees use signs meaningfully but lack syntax. </li></ul><ul><li>Presented with ambiguous information, people tend to see what they want to see. </li></ul> Conclusions <ul><li>If we say that animals can use meaningful sequences of signs to communicate a capability for language, our understanding would be naive… Steven Pinker (1995) concludes, “chimps do not develop language.” </li></ul> A particular slide catching your eye? Clipping is a handy way to collect important slides you want to go back to later.
Memories in Circles Match the different terms related to circles to its description. Around the Circle Identify the different parts of a circle along with the terms associated with it. Gill's Angles: Angles, Triangles, Parallel Lines, Transversals Angles, Triangles, Parallel Lines, Transversals Geometry review questions Pythagorean Theorem: Midpoint, Perimeter and Area Unit 5 Lesson 2: We will work on finding the midpoint, the Pythagorean theorem and finding the area and perimeter of triangles and rectangles. Angles: Match Three Types of Angles Match three types of angles to their definitions Using your memory, quickly match the correct shapes! Supplementary Angle Match Students will choose cards that are supplementary to each other. Review different types of polygons while playing a fun game!
- The scope of astronomy - Determining astronomical distances - Study of the solar system - Study of the stars - Study of the Milky Way Galaxy - Study of other galaxies and related phenomena - The techniques of astronomy - Impact of astronomy - History of astronomy - Prehistory and antiquity - India, the Islamic world, medieval Europe, and China - The age of observation - The rise of astrophysics - Galaxies and the expanding universe - The origin of the universe - Echoes of the big bang The motion of the planets Greek thinking about the motion of the planets began by about 400 bce. Eudoxus of Cnidus constructed the first Greek theory of planetary motion of which any details are known. In a book, On Speeds (which is lost but was briefly discussed by Aristotle and Simplicius), Eudoxus regarded each celestial body as carried on a set of concentric spheres, which nest one inside another. For each planet, three different motions must be accounted for, and Eudoxus proposed to do this with four spheres. The daily revolution to the west is accounted for by the outermost sphere (1). Next inside is sphere 2, whose axle fits into sphere 1 at an offset of about 24°; sphere 2 turns to the east in the planet’s zodiacal period (12 years for Jupiter, 30 years for Saturn). The third motion is retrograde motion. For this, Eudoxus used a combination of two spheres (3 and 4). The planet itself rides on the equator circle of sphere 4. The axle of 4 fits inside sphere 3 with a slight angular offset. Spheres 3 and 4 turn in opposite directions but at the same speed. The motion of the planet resulting from the gyrations of spheres 3 and 4 is a figure eight, which lies in the spherical surface. Eudoxus likely understood the mathematical characteristics of this curve, as he gave it the name hippopede (horse fetter). The two-sphere assembly of 3 and 4 is inserted into the inner surface of sphere 2. Thus, all three motions are accounted for, at least qualitatively: the daily motion to the west by sphere 1, the slow motion eastward around the zodiac by sphere 2, and the occasional retrograde motion by the two-sphere assembly of 3 and 4. Eudoxus’s theory is sometimes called the theory of homocentric spheres, as all the spheres have the same centre, Earth. At this stage, Greek astronomers were more interested in providing plausible physical accounts of the universe and in proving geometrical theorems than in providing numerically accurate descriptions of planetary motion. Eudoxus’s successor Callippus made some improvements to the model. Nevertheless, the homocentric spheres were criticized for their failure to account for the fact that some planets (notably Mars and Venus) are much brighter at some times of their cycles than at others. Eudoxus’s system was soon abandoned as a theory for the motion of the planets, but it exerted a profound influence in cosmology, for the cosmos continued to be regarded as a set of concentric spheres until the Renaissance. Late in the 3rd century bce, alternative theoretical models were developed, based on eccentric circles and epicycles. (An eccentric circle is a circle that is slightly off-centre from Earth, and an epicycle is a circle that is carried and rides around on another circle.) This innovation is usually attributed to Apollonius of Perga (c. 220 bce), but it is not conclusively known who first proposed these models. In considering the Sun’s motion, Eudoxus’s theory of homocentric spheres ignored the fact that the Sun appears to speed up and slow down in the course of the year as it moves around the zodiac. (This is clear from spring’s being several days longer than fall.) An eccentric (i.e., off-centre) circle can explain this fact. The Sun is still considered to travel at constant speed around a perfect circle, but the centre of the circle is slightly displaced from Earth. When the Sun is closest to Earth, it appears to travel a little more rapidly in the zodiac. When it is farthest away, it appears to travel a little more slowly. As far as is known, Hipparchus was the first to deduce the amount and direction of the off-centredness, basing his calculations on the measured length of the seasons. According to Hipparchus, the off-centredness of the Sun’s circle is about 4 percent of its radius. The eccentric-circle theory was capable of excellent accuracy in accounting for the observed motion of the Sun and remained standard until the 17th century. The standard theory of the planets involved an eccentric circle, which carried an epicycle. Imagine looking down on the plane of the solar system from above its north pole. The planet moves counterclockwise on its epicycle. Meanwhile, the centre of the epicycle moves counterclockwise around the eccentric circle, which is centred near (but not quite exactly at) Earth. As viewed from Earth, the planet will appear to move backward (that is, go into retrograde motion) when it is at the inner part of the epicycle (closest to Earth), for this is when the westward motion of the planet on the epicycle is more than enough to overcome the eastward motion of the epicycle’s centre forward around the eccentric. Hipparchus played a major role in introducing Babylonian numerical parameters into Greek astronomy. Indeed, an important shift in Greek attitudes toward astronomy occurred about this time. The Babylonian example served as a sort of wake-up call to the Greeks. Previous Greek planetary thinking had been more about getting the right big picture, based on philosophical principles and geometrical models (whether using Eudoxus’s concentric spheres or Apollonius’s epicycles and eccentrics). The Babylonians had no geometrical models but instead focused on devising arithmetical theories that had real predictive power. Hipparchus achieved numerically successful geometrical theories for the Sun and the Moon, but he did not succeed with the planets. He contented himself with showing that the planetary theories then in circulation did not agree with the phenomena. Nevertheless, Hipparchus’s insistence that a geometrical theory, if it is true, ought to work in detail marked a major step in Greek astronomy. Another of Hipparchus’s contributions was the discovery of precession, the slow eastward motion of the stars around the zodiac caused by wobbling, over a period of 25,772 years, in the orientation of Earth’s axis of rotation. Hipparchus’s writings on this subject have not survived, but his ideas can be reconstructed from summaries given by Ptolemy. Hipparchus used observations of several fixed stars, taken with respect to the eclipsed Moon, which had been made by some of his predecessors. On comparing these with eclipse observations he had made himself, he deduced that the fixed stars move eastward not less than 1° in 100 years. The Babylonians, in their theories, revised their locations of the equinoxes and solstices. For example, in one version of the Babylonian theory, the spring equinox is said to occur at the 10th degree of Aries; in another version, at the 8th degree. Some historians have maintained that this reflects a Babylonian awareness of precession, on which Hipparchus might have drawn. Other historians have argued that the evidence is not clear and that these differing norms for the equinox may represent nothing more than alternative conventions.
Learn the definition of chemical equilibrium and how it is dynamic. Discover what the equilibrium constant is and how it shows whether the reaction favors the reactants or products. Learn how chemists designate equilibrium in an equation and how they show the difference in reaction rate. Chemical equilibrium is when the rate of a forward reaction equals the rate of the reverse reaction and the concentrations of the products and reactants remain unchanged. Equilibrium is a dynamic state, meaning that things are always moving. Products are being broken down into reactants, and reactants are being combined into products. Things are moving, but the concentrations stay the same. When the reaction is written, it is written with a double arrow instead of an equal sign to show that the reaction is reversible. A + B go to C + D. In some reactions, the forward reaction is almost completed before the reverse reaction starts. In this case, there's a higher concentration of products than reactants, but the reaction can still be in equilibrium because the concentrations of both the reactants and products stay the same. The reaction equilibrium lies to the right because there are more products than reactants. In this case, the reaction is written with two different-length arrows, with the longer arrow pointing to the right, showing that more product is made than reactant. A + B goes to C + D. The opposite is also true. The forward reaction of making products has barely started, and the reverse reaction is already going like gangbusters. In this case, the equilibrium of the reaction is said to lie to the left and the longer arrow points left. A + B goes to C + D. The rate of reactions are often shown in a graph like this one. This graph compares the rate of the forward reaction to the rate of the reverse reaction. To start, the forward reaction has the maximum rate possible, and the reverse reaction has no rate because it hasn't started yet. As the reaction is under way, the forward reaction decreases as the reactants are used up, and the reverse reaction increases as there is more product to turn back into reactants. Eventually, equilibrium is reached, and the graph turns into one straight horizontal line. Once equilibrium is reached, the concentrations of the reactants and products don't change. When this happens, an equilibrium constant K can be written for the reaction. Only the substances whose concentrations change are included in the equilibrium constant equation. If the reaction equation is nA + mB goes to xC + yD, then the equilibrium constant equation can be written as: The brackets indicate concentration. The superscripts are the same as the coefficients of each substance in the balanced equation. Let me show you how this looks in a real balanced equation and it may make it clearer for you. The K for a reaction at a given temperature shows how many reactants are converted into product. If the K is small, the forward reaction barely starts before the reverse reaction gets going and equilibrium is established. If K is large, then most of the reactants were made into products before the reaction reaches equilibrium. You can solve for any of the unknowns in an equilibrium constant equation if you know the other information. Example: Determine the K for the reaction N2 + O2 goes to 2NO. - [N2] = 0.0064 mol/L - [O2] = 0.0017 mol/L - [NO] = 1.1 * 10-5 mol/L Here's another example: Determine the [I2] in the following equation: H2 + I2 goes to 2HI. - K = 54 - [HI] = 0.017 - [H2] = 0.002 Chemical equilibrium is when the rate of a forward reaction equals the rate of the reverse reaction and the concentrations of the products and reactants remain unchanged. Equilibrium is a dynamic state, meaning that things are always moving. Once equilibrium is reached, the concentrations of the reactants and products don't change, and an equilibrium constant equation can be written for the reaction. If the reaction equation is nA + mB goes to xC + yD, then the equilibrium constant equation can be written as: After viewing this video, you should be able to paraphrase what chemical equilibrium is and how it is always in motion but has equal reactions within it.
3Circular MotionWhen an object moves in a circle at constant speed, we describe it as undergoing uniform circular motion.Its speed is constant, but its velocity is not because velocity includes direction and the object’s direction is clearly changing. 4Circular Motion A changing velocity means acceleration. The pull on the string is always directed perpendicular to the velocity.The pull accelerates the ball into a circular path, even though the ball does not speed up or slow down.The pull changes only the direction of the velocity, not the magnitude. 8Centripetal Acceleration Centripetal means center-seeking.Centripetal acceleration is always directed toward the center of the circle of motion.It is this centripetal acceleration that is responsible for the change in the direction of the velocity; the magnitude of the velocity remains constant.Any change in the tangential acceleration causes a change in the speed of the particle as it travels around the circle. In uniform circular motion, aT = 0, so the acceleration is completely radial (ar) or centripetal. 9Centripetal ForceNewton’s Second Law explains that an object undergoing acceleration is experiencing a net force. The net force on an object undergoing uniform circular motion is called the centripetal force Fc.The centripetal force necessary for an object of mass m to travel with constant speed v in a circle of radius r is given by: 10Centripetal ForceThe centripetal force always points toward the center of the circle about which the object moves with uniform speed.If the centripetal force applied to the object is removed, the object will move in a straight-line tangent to the curved path at the point where the centripetal force ceases. When the centripetal force ceases, the object has no unbalanced forces acting upon it and thus moves in a straight line at constant speed. 11Circular MotionIf the string breaks, the ball flies off in a straight line. It is the force of the string that causes the acceleration in this example of uniform circular motion. 12Centripetal ForceCentripetal force is the name given to any force that is directed at right angles to the path of a moving object and that tends to produce circular motion.Examples:the gravitational force directed toward the center of the Earth holds the Moon in an almost circular orbit about the Earth;an electromagnetic force that is directed toward the nucleus holds the electrons that revolve about the nucleus of the atom.Directions in centripetal force problems:Positive direction is inwards toward center of circle.Negative direction is outward away from center of circle. 13Radius r is the distance from the center of the mass to the axis of rotation. 14The centripetal force is not a force and does not belong in a free-body diagram.The force F in the picture would provide the centripetal force needed to maintain the circular path. 15Motion On A Flat CurveThe net force on a car traveling around a curve is the centripetal force.As a car travels around a curve, the net force on the car must be the centripetal force, directed toward the center of the circle the curve is a portion of. 16Motion On A Flat CurveOn a flat, level curve, the friction between the tires and the road supplies the centripetal force.If the tires are worn smooth or the road is icy or oily, this friction force will not be available.The car will not be able to move in a circle, it will keep going in a straight line and therefore go off the road. 17Motion On A Flat Curve Accelerations: Equation: Fc = FF ; ay = 0 m/s2 ax = ac =Equation: Fc = FF ; 18Motion On A Banked Curve Some curves are banked to compensate for slippery conditions.In addition to any friction forces that may or may not be present, the road exerts a normal force perpendicular to its surface.The downward force of the car’s weight is also present.These two forces add as vectors to provide a net force Fnet that points toward the center of the circle; this is the centripetal force.The centripetal force is directed toward the center of the circle, it is not parallel to the banked road. 19Motion On A Banked Curve The effect of banking is to tilt the normal force Fn toward the center of curvature of the road so that the inward radial component FNsin can supply the required centripetal force.Vehicles can make a sharp turn more safely if the road is banked. If the vehicle maintains the speed for which the curve is designed, no frictional force is needed to keep the vehicle on the road. 20Motion On A Banked Curve The effect of banking is to tilt the normal force FN toward the center of curvature of the road so that the inward radial component FNsin can supply the required centripetal force. 21Motion On A Banked Curve There is no acceleration along the y axis, so the sum of the forces in the y plane is zero:The horizontal component of the normal force FN, the force the road exerts against the car, provides the necessary centripetal force. Because the only force in the x plane is the centripetal force: 22Motion On A Banked Curve This equation gives the banking angle that allows a car to travel in a curve of radius r with constant speed v and require no friction force.Goldilocks: Just Right! 23Motion On A Banked Curve A banked curve is designed for one specific speed, called the “design speed”.If the banked curve is icy so that there is no friction force at all, then traveling at a speed higher than the design speed means the car will slide out, up, and over the edge.Traveling at a lower speed than the design speed means the car will slide in, down, and off the bank.When a banked curve has to be negotiated at a speed above the design speed (or if you are asked to find the maximum speed), friction is required. The frictional force Ff acts parallel to the road surface. 24Motion On A Banked Curve (w/ Friction) Goldilocks: Too Fast! 25Motion On A Banked Curve (w/ Friction) Goldilocks: Too Fast! 26Motion On A Banked Curve (w/ Friction) Resolve the normal force FN and friction FF into horizontal and vertical components.The horizontal components of FN and FF are both directed inward toward the center of the curve, therefore, these two force components combine to determine the centripetal force; ΣFx = m·ac. 27Motion On A Banked Curve (w/ Friction) Because there are no unbalanced forces in the vertical direction, the upward forces must equal the downward forces, therefore:Substitute into both equations. 28Motion On A Banked Curve (w/ Friction) Solve both equations for Fn and set these equations equal to each other (because FN = FN as there is only one force acting normal to the surface). From this equation, the unknown variable can be determined. 29Motion On A Banked Curve (w/ Friction) For problems involving a minimum speed for the vehicle to travel around the curve without skidding, the frictional force is directed up the incline to keep the vehicle from sliding down to the bottom of the banked curve.Resolve the normal force FN and friction FF into horizontal and vertical components. 30Motion On A Banked Curve (w/ Friction) Goldilocks: Too Slow! 31Motion On A Banked Curve (w/ Friction) Goldilocks: Too Slow! 32Motion On A Banked Curve (w/ Friction) The horizontal component of FN is inward toward the center of the curve and is positive; FF is directed outward away from the center of the curve and is negative. These two force components combine to determine the centripetal force; ΣFx = m·ac. 33Because there are no unbalanced forces in the vertical direction, the upward forces must equal the downward forces, therefore:Substitute into both equations. 34Solve both equations for Fn and set these equations equal to each other (because FN = FN as there is only one force acting normal to the surface). From this equation, the unknown variable can be determined. 35Vertical CirclesThe force of gravity causes the speed of an object in a vertical circular path to vary. The object accelerates on the downward portion of its circular path and decelerates on the upward portion of the circular path.At the top and bottom of a vertical circular path, the weight and the normal force (or an equivalent supporting force, such as tension) are the only forces acting on an object. The centripetal force is supplied by the resultant of the weight and a supporting force (often the normal force). 36Vertical CirclesThe forces acting on a person sitting in a roller coaster car are shown. The person’s weight FW is present and so is the normal force FN that the seat exerts on him (this is your apparent weight). 37Vertical CirclesThe normal force FN, the force you feel on the seat of your pants, can be positive, negative, or zero.A negative value for FN means the passenger has to be strapped in, with the straps exerting an upward force. Such a situation would be dangerous, and roller coaster designers avoid this.If FN = 0 N, the person seems to be weightless as well as upside down. 38Vertical CirclesThe forces on an airplane pilot at the bottom of a dive can be quite large.Gravity pulls downward and the seat exerts its usual normal force FN, this time upward. 39Vertical CirclesAt the bottom of the dive, the normal force can only be positive, must be greater than the weight, and can become very large. A roller coaster at the bottom of the loop provides the same forces.The acceleration can be expressed as:This acceleration can be expressed in terms of g’s, where g’s are determined by dividing the centripetal acceleration by gravity. One g is m/s2. The number of g’s represents the relative pull of gravity on the body that the person experiences. 40Vertical CirclesExperiencing a significant number of g’s makes the work of the heart more difficult. Accelerations of eight to ten g’s make it difficult for the circulatory system to get enough blood to the brain and may result in blackouts. Pressure suits that squeeze on the legs push blood back into the rest of the body, including the brain, and help prevent blackouts. 41Vertical CirclesFor the Ferris wheel, the only difference occurs at the top where the seat is facing upward.Top:This equation is also for a car passing over the top of a curve.Bottom: 42TensionFor an object attached to a string and moving in a vertical circle, the centripetal force is at a minimum at the top of its vertical path and at a maximum at the bottom of its vertical path. 43TensionTop: the centripetal force on the object equals the tension of the string plus the weight of the ball, both acting toward the center of the vertical circle. Mathematically:Bottom: the centripetal force on the object is equal to the difference between the tension of the string and the weight of the object. The tension is exerted inward toward the center of the vertical circle, while the weight is directed away from the center of the vertical circle. Mathematically: 44Critical Velocity (vmin) Critical velocity: velocity below which an object moving in a vertical circle will not describe a circular path.Critical velocity depends on the acceleration due to gravity and the radius of the vertical circle, not on the mass of the object. 45Tension at an Angle & Horizontal Circles Vertical component of the tension is equal to the weight.Horizontal component of the weight is equal to the centripetal force.the radius is the distance from the center of the mass to the dot at the center of the horizontal circle. 46Conical Pendulum L T r Fw = m·g For conical pendulums, centripetal force is provided by a component of the tension. is the angle between the vertical and the cord.TFw = m·gisLr 49Fn Difference?In every previous inclined plane problem such as a skier on a hill or a block on a ramp, the normal force is given by Fn = m·g·cos .In every banked curve problem the normal force is given by Fn=m·g/cos .Both equations are correct!The difference is due to the direction of the acceleration.the direction of the acceleration in the inclined plane problem is down toward the bottom of the incline. There is no component of the acceleration in the normal direction. 50the direction of the acceleration in the banked curve is horizontally toward the center of the circle. This produces a component of acceleration in the normal direction.In the inclined plane problems, I taught you to rotate the x- and y- axes so that the x-axis is parallel to the surface and the y-axis is perpendicular to the surface. We did not do this with the banked turn problems.This results in the acceleration being zero on the y axis for both cases.In the inclined plane problems, the acceleration is zero along the normal perpendicular to the incline; thus the normal force equals a component of gravity.) 51In the frictionless banked turn, the acceleration is zero vertically; thus the force of gravity equals a component of the normal force.In the inclined plane problems, one component of gravity causes the acceleration (the other component cancels out the normal force).In the banked turn, one component of the normal force causes the acceleration (the other component cancels out the force of gravity). 52For Uniform Circular Motion: n = number of revolutions (rotations)r = radiust = timePeriod = time for one revolution (rotation) 53Web Sites Amusement Park Physics Roller Coasters and Amusement Park PhysicsCoasterQuest.com
Thermodynamic temperature is a quantity defined in thermodynamics as distinct from kinetic theory or statistical mechanics. A thermodynamic temperature reading of zero is of particular importance for the third law of thermodynamics. By convention, it is reported on the Kelvin scale of temperature in which the unit of measurement is the kelvin (unit symbol: K). For comparison, a temperature of 295 K is equal to 21.85 °C and 71.33 °F. At the zero point of thermodynamic temperature, absolute zero, the particle constituents of matter have minimal motion and can become no colder. Absolute zero, which is a temperature of zero kelvins (0 K), is precisely equal to −273.15 °C and −459.67 °F. Matter at absolute zero has no remaining transferable average kinetic energy and the only remaining particle motion is due to an ever-pervasive quantum mechanical phenomenon called zero-point energy. Though the atoms in, for instance, a container of liquid helium that was precisely at absolute zero would still jostle slightly due to zero-point energy, a theoretically perfect heat engine with such helium as one of its working fluids could never transfer any net kinetic energy (heat energy) to the other working fluid and no thermodynamic work could occur. Temperature is generally expressed in absolute terms when scientifically examining temperature's interrelationships with certain other physical properties of matter such as its volume or pressure (see Gay-Lussac's law), or the wavelength of its emitted black-body radiation. Absolute temperature is also useful when calculating chemical reaction rates (see Arrhenius equation). Furthermore, absolute temperature is typically used in cryogenics and related phenomena like superconductivity, as per the following example usage: “Conveniently, tantalum’s transition temperature (Tc) of 4.4924 kelvin is slightly above the 4.2221 K boiling point of helium.” The International System of Units (SI) specifies the Kelvin scale for measuring thermodynamic temperature, and the unit of measure kelvin (unit symbol: K) for specific values along the scale. The kelvin is also used for denoting temperature intervals (a span or difference between two temperatures) as per the following example usage: “A 60/40 tin/lead solder is non-eutectic and is plastic through a range of 5 kelvins as it solidifies.” A temperature interval of one degree Celsius is the same magnitude as one kelvin. The magnitude of the kelvin was redefined in 2019 in relation to the very physical property underlying thermodynamic temperature: the kinetic energy of atomic particle motion. The redefinition fixed the Boltzmann constant at precisely 1.380649×10−23 joules per kelvin (J/K). The compound unit of measure for the Boltzmann constant is often also given as J·K−1, which may seem abstract due the multiplication dot (·) and a kelvin symbol that is followed by a superscripted negative 1 exponent. However, this is merely another mathematical syntax denoting the same measure: joules (the SI unit for energy, including kinetic energy) per kelvin. The property that imbues any substances with a temperature can be readily understood by examining the ideal gas law, which relates, per the Boltzmann constant, how heat energy causes precisely defined changes in the pressure and temperature of certain gases. This is because monatomic gases like helium and argon behave kinetically like perfectly elastic and spherical billiard balls that move only in a specific subset of the possible motions that can occur in matter: that comprising the three translational degrees of freedom. The translational degrees of freedom are the familiar billiard ball-like movements along the X, Y, and Z axes of 3D space (see Fig. 1, below). This is why the noble gases all have the same specific heat capacity per atom and why that value is lowest of all the gases. Molecules (two or more chemically bound atoms), however, have internal structure and therefore have additional internal degrees of freedom (see Fig. 3, below), which makes molecules absorb more heat energy for any given amount of temperature rise than do the monatomic gases. Heat energy is born in all available degrees of freedom; this is in accordance with the equipartition theorem, so all available internal degrees of freedom have the same temperature as their three external degrees of freedom. However, the property that gives all gases their pressure, which is the net force per unit area on a container arising from gas particles recoiling off it, is a function of the kinetic energy borne in the atoms’ and molecules’ three translational degrees of freedom. Fixing the Boltzmann constant at specific value, along with other rule making, had the effect of precisely establishing the magnitude of the unit interval of thermodynamic temperature, the kelvin, in terms of the average kinetic behavior of the noble gases. Moreover, the starting point of the thermodynamic temperature scale, absolute zero, was reaffirmed as the point at which zero average kinetic energy remains in a sample; the only remaining particle motion being that comprising random vibrations due to zero-point energy. The Rankine scaleEdit Though there have been many other temperature scales throughout history, there have been only two scales for measuring thermodynamic temperature where absolute zero is their null point (0): The Kelvin scale and the Rankine scale. Throughout the scientific world where modern measurements are nearly always made using the International System of Units, thermodynamic temperature is measured using the Kelvin scale. The Rankine scale is part of English engineering units in the United States and finds use in certain engineering fields, particularly in legacy reference works. The Rankine scale uses the degree Rankine (symbol: °R) as its unit, which is the same magnitude as the degree Fahrenheit (symbol: °F). A unit increment of one degree Rankine is precisely 1.8 times smaller in magnitude than one kelvin; thus, to convert a specific temperature on the Kelvin scale to the Rankine scale, K × 1.8 = °R, and to convert from a temperature on the Rankine scale to the Kelvin scale, °R / 1.8 = K. Consequently, absolute zero is “0” for both scales, but the melting point of water ice (0 °C and 273.15 K) is 491.67 °R. To convert temperature intervals (a span or difference between two temperatures), one uses the same formulas from the preceding paragraph; for instance, a range of 5 kelvins is precisely equal to a range of 9 degrees Rankine. Modern redefinition of the kelvinEdit For 65 years, between 1954 and the 2019 redefinition of the SI base units, a temperature interval of one kelvin was defined as 1/273.16 the difference between the triple point of water and absolute zero. The 1954 resolution by the International Bureau of Weights and Measures (known by the French-language acronym BIPM), plus later resolutions and publications, defined the triple point of water as precisely 273.16 K and acknowledged that it was “common practice” to accept that due to previous conventions (namely, that 0 °C had long been defined as the melting point of water and that the triple point of water had long been experimentally determined to be indistinguishably close to 0.01 °C), the difference between the Celsius scale and Kelvin scale is accepted as 273.15 kelvins; which is to say, 0 °C equals 273.15 kelvins. The net effect of this as well as later resolutions was twofold: 1) they defined absolute zero as precisely 0 K, and 2) they defined that the triple point of special isotopically controlled water called Vienna Standard Mean Ocean Water was precisely 273.16 kelvins and 0.01 °C. One effect of the aforementioned resolutions was that the melting point of water, while very close to 273.15 kelvin and 0 °C, was not a defining value and was subject to refinement with more precise measurements. The 1954 BIPM standard did a good job of establishing—within the uncertainties due to isotopic variations between water samples—temperatures around the freezing and triple points of water, but required that intermediate values between the triple point and absolute zero, as well as extrapolated values from room temperature and beyond, to be experimentally determined via apparatus and procedures in individual labs. This shortcoming was addressed by the International Temperature Scale of 1990, or ITS‑90, which defined 13 additional points, from 13.8033 K, to 1,357.77 K. While definitional, ITS‑90 had—and still has—some challenges, partly because eight of its extrapolated values depend upon the melting or freezing points of metal samples, which must remain exceedingly pure lest their melting or freezing points be affected—usually depressed. The 2019 redefinition of the SI base units was primarily for the purpose of decoupling much of the SI system's definitional underpinnings from the kilogram, which was the last physical artifact defining an SI base unit (a platinum/iridium cylinder stored under three nested bell jars in a safe located in France) and which had highly questionable stability. The solution required that four physical constants, including the Boltzmann constant, be definitionally fixed. Assigning the Boltzmann constant a precisely defined value had no practical effect on modern thermometry except for the most exquisitely precise measurements. Before the redefinition, the triple point of water was exactly 273.16 K and 0.01 °C and the Boltzmann constant was experimentally determined to be 1.38064903(51)×10−23 J/K, where the “(51)” denotes the uncertainty in the two least significant digits (the 03) and equals a relative standard uncertainty of 0.37 ppm. Afterwards, by defining the Boltzmann constant as exactly 1.380649×10−23 J/K, the 0.37 ppm uncertainty was transferred to the triple point of water, which became an experimentally determined value of 273.1600 ±0.0001 K (0.0100 ±0.0001 °C). That the triple point of water ended up being exceedingly close to 273.16 K after the SI redefinition was no accident; the final value of the Boltzmann constant was determined, in part, through clever experiments with argon and helium that used the triple point of water for their key reference temperature. Notwithstanding the 2019 redefinition, water triple-point cells continue to serve in modern thermometry as exceedingly precise calibration references at 273.16 K and 0.01 °C. Moreover, the triple point of water remains one of the 14 calibration points comprising ITS‑90, which spans from the triple point of hydrogen (13.8033 K) to the freezing point of copper (1,357.77 K), which is a nearly hundredfold range of thermodynamic temperature. The relationship of temperature, motions, conduction, and thermal energyEdit The nature of kinetic energy, translational motion, and temperatureEdit The thermodynamic temperature of any bulk quantity of a substance (a statistically significant quantity of particles) is directly proportional to the mean average kinetic energy of a specific kind of particle motion known as translational motion. These simple movements in the three X, Y, and Z–axis dimensions of space means the particles move in the three spatial degrees of freedom. This particular form of kinetic energy is sometimes referred to as kinetic temperature. Translational motion is but one form of heat energy and is what gives gases not only their temperature, but also their pressure and the vast majority of their volume. This relationship between the temperature, pressure, and volume of gases is established by the ideal gas law's formula pV = nRT and is embodied in the gas laws. Though the kinetic energy borne exclusively in the three translational degrees of freedom comprise the thermodynamic temperature of a substance, molecules, as can be seen in Fig. 3, can have other degrees of freedom, all of which fall under three categories: bond length, bond angle, and rotational. All three additional categories are not necessarily available to all molecules, and even for molecules that can experience all three, some can be “frozen out” below a certain temperature. Nonetheless, all those degrees of freedom that are available to the molecules under a particular set of conditions contribute to the specific heat capacity of a substance; which is to say, they increase the amount of heat (kinetic energy) required to raise a given amount of the substance by one kelvin or one degree Celsius. The relationship of kinetic energy, mass, and velocity is given by the formula Ek = 1/2mv2. Accordingly, particles with one unit of mass moving at one unit of velocity have precisely the same kinetic energy, and precisely the same temperature, as those with four times the mass but half the velocity. The extent to which the kinetic energy of translational motion in a statistically significant collection of atoms or molecules in a gas contributes to the pressure and volume of that gas is a proportional function of thermodynamic temperature as established by the Boltzmann constant (symbol: kB). The Boltzmann constant also relates the thermodynamic temperature of a gas to the mean kinetic energy of an individual particles’ translational motion as follows: - is the mean kinetic energy for any individual particle, in joules (J) - kB = 1.380649×10−23 J/K - T is the thermodynamic temperature of the bulk quantity of the substance, in kelvins (K) While the Boltzmann constant is useful for finding the mean kinetic energy in a sample of particles, it's important to note that even when a substance is isolated and in thermodynamic equilibrium (all parts are at a uniform temperature and no heat is going into or out of it), the translational motions of individual atoms and molecules occurs across a wide range of speeds (see animation in Fig. 1 above). At any one instant, the proportion of particles moving at a given speed within this range is determined by probability as described by the Maxwell–Boltzmann distribution. The graph shown here in Fig. 2 shows the speed distribution of 5500 K helium atoms. They have a most probable speed of 4.780 km/s (0.2092 s/km). However, a certain proportion of atoms at any given instant are moving faster while others are moving relatively slowly; some are momentarily at a virtual standstill (off the x–axis to the right). This graph uses inverse speed for its x–axis so the shape of the curve can easily be compared to the curves in Fig. 5 below. In both graphs, zero on the x–axis represents infinite temperature. Additionally, the x and y–axis on both graphs are scaled proportionally. The high speeds of translational motionEdit Although very specialized laboratory equipment is required to directly detect translational motions, the resultant collisions by atoms or molecules with small particles suspended in a fluid produces Brownian motion that can be seen with an ordinary microscope. The translational motions of elementary particles are very fast and temperatures close to absolute zero are required to directly observe them. For instance, when scientists at the NIST achieved a record-setting cold temperature of 700 nK (billionths of a kelvin) in 1994, they used optical lattice laser equipment to adiabatically cool cesium atoms. They then turned off the entrapment lasers and directly measured atom velocities of 7 mm per second to in order to calculate their temperature. Formulas for calculating the velocity and speed of translational motion are given in the following footnote. It is neither difficult to imagine atomic motions due to kinetic temperature, nor distinguish between such motions and those due to zero-point energy. Consider the following hypothetical thought experiment, as illustrated in Fig. 2.5 at left, with an atom that is exceedingly close to absolute zero. Imagine peering through a common optical microscope set to 400 power, which is about the maximum practical magnification for optical microscopes. Such microscopes generally provide fields of view a bit over 0.4 mm in diameter. At the center of the field of view is a single levitated argon atom (argon comprises about 0.93% of air) that is illuminated and glowing against a dark backdrop. If this argon atom was at a beyond-record-setting one-trillionth of a kelvin above absolute zero, and was moving perpendicular to the field of view towards the right, it would require 13.9 seconds to move from the center of the image to the 200-micron tick mark; this travel distance is about the same as the width of the period at the end of this sentence on modern computer monitors. As the argon atom slowly moved, the positional jitter due to zero-point energy would be much less than the 200-nanometer (0.0002 mm) resolution of an optical microscope. Importantly, the atom's translational velocity of 14.43 microns per second constitutes all its retained kinetic energy due to not being precisely at absolute zero. Were the atom precisely at absolute zero, imperceptible jostling due to zero-point energy would cause it to very slightly wander, but the atom would perpetually be located, on average, at the same spot within the field of view. This is analogous to a boat that has had its motor turned off and is now bobbing slightly in relatively calm and windless ocean waters; even though the boat randomly drifts to and fro, it stays in the same spot in the long term and makes no headway through the water. Accordingly, an atom that was precisely at absolute zero would not be “motionless,” and yet, a statistically significant collection of such atoms would have zero net kinetic energy available to transfer to any other collection of atoms. This is because regardless of the kinetic temperature of the second collection of atoms, they too experience the effects of zero-point energy. Such are the consequences of statistical mechanics and the nature of thermodynamics. The internal motions of molecules and internal energyEdit As mentioned above, there are other ways molecules can jiggle besides the three translational degrees of freedom that imbue substances with their kinetic temperature. As can be seen in the animation at right, molecules are complex objects; they are a population of atoms and thermal agitation can strain their internal chemical bonds in three different ways: via rotation, bond length, and bond angle movements; these are all types of internal degrees of freedom. This makes molecules distinct from monatomic substances (consisting of individual atoms) like the noble gases helium and argon, which have only the three translational degrees of freedom (the X, Y, and Z axis). Kinetic energy is stored in molecules’ internal degrees of freedom, which gives them an internal temperature. Even though these motions are called “internal,” the external portions of molecules still move—rather like the jiggling of a stationary water balloon. This permits the two-way exchange of kinetic energy between internal motions and translational motions with each molecular collision. Accordingly, as internal energy is removed from molecules, both their kinetic temperature (the kinetic energy of translational motion) and their internal temperature simultaneously diminish in equal proportions. This phenomenon is described by the equipartition theorem, which states that for any bulk quantity of a substance in equilibrium, the kinetic energy of particle motion is evenly distributed among all the active degrees of freedom available to the particles. Since the internal temperature of molecules are usually equal to their kinetic temperature, the distinction is usually of interest only in the detailed study of non-local thermodynamic equilibrium (LTE) phenomena such as combustion, the sublimation of solids, and the diffusion of hot gases in a partial vacuum. The kinetic energy stored internally in molecules causes substances to contain more heat energy at any given temperature and to absorb additional internal energy for a given temperature increase. This is because any kinetic energy that is, at a given instant, bound in internal motions is not at that same instant contributing to the molecules’ translational motions. This extra kinetic energy simply increases the amount of internal energy a substance absorbs for a given temperature rise. This property is known as a substance's specific heat capacity. Different molecules absorb different amounts of internal energy for each incremental increase in temperature; that is, they have different specific heat capacities. High specific heat capacity arises, in part, because certain substances’ molecules possess more internal degrees of freedom than others do. For instance, room-temperature nitrogen, which is a diatomic molecule, has five active degrees of freedom: the three comprising translational motion plus two rotational degrees of freedom internally. Not surprisingly, in accordance with the equipartition theorem, nitrogen has five-thirds the specific heat capacity per mole (a specific number of molecules) as do the monatomic gases. Another example is gasoline (see table showing its specific heat capacity). Gasoline can absorb a large amount of heat energy per mole with only a modest temperature change because each molecule comprises an average of 21 atoms and therefore has many internal degrees of freedom. Even larger, more complex molecules can have dozens of internal degrees of freedom. The diffusion of thermal energy: Entropy, phonons, and mobile conduction electronsEdit Heat conduction is the diffusion of thermal energy from hot parts of a system to cold parts. A system can be either a single bulk entity or a plurality of discrete bulk entities. The term bulk in this context means a statistically significant quantity of particles (which can be a microscopic amount). Whenever thermal energy diffuses within an isolated system, temperature differences within the system decrease (and entropy increases). One particular heat conduction mechanism occurs when translational motion, the particle motion underlying temperature, transfers momentum from particle to particle in collisions. In gases, these translational motions are of the nature shown above in Fig. 1. As can be seen in that animation, not only does momentum (heat) diffuse throughout the volume of the gas through serial collisions, but entire molecules or atoms can move forward into new territory, bringing their kinetic energy with them. Consequently, temperature differences equalize throughout gases very quickly—especially for light atoms or molecules; convection speeds this process even more. Translational motion in solids, however, takes the form of phonons (see Fig. 4 at right). Phonons are constrained, quantized wave packets that travel at the speed of sound of a given substance. The manner in which phonons interact within a solid determines a variety of its properties, including its thermal conductivity. In electrically insulating solids, phonon-based heat conduction is usually inefficient and such solids are considered thermal insulators (such as glass, plastic, rubber, ceramic, and rock). This is because in solids, atoms and molecules are locked into place relative to their neighbors and are not free to roam. Metals however, are not restricted to only phonon-based heat conduction. Thermal energy conducts through metals extraordinarily quickly because instead of direct molecule-to-molecule collisions, the vast majority of thermal energy is mediated via very light, mobile conduction electrons. This is why there is a near-perfect correlation between metals' thermal conductivity and their electrical conductivity. Conduction electrons imbue metals with their extraordinary conductivity because they are delocalized (i.e., not tied to a specific atom) and behave rather like a sort of quantum gas due to the effects of zero-point energy (for more on ZPE, see Note 1 below). Furthermore, electrons are relatively light with a rest mass only 1⁄1836 that of a proton. This is about the same ratio as a .22 Short bullet (29 grains or 1.88 g) compared to the rifle that shoots it. As Isaac Newton wrote with his third law of motion, Law #3: All forces occur in pairs, and these two forces are equal in magnitude and opposite in direction. However, a bullet accelerates faster than a rifle given an equal force. Since kinetic energy increases as the square of velocity, nearly all the kinetic energy goes into the bullet, not the rifle, even though both experience the same force from the expanding propellant gases. In the same manner, because they are much less massive, thermal energy is readily borne by mobile conduction electrons. Additionally, because they're delocalized and very fast, kinetic thermal energy conducts extremely quickly through metals with abundant conduction electrons. The diffusion of thermal energy: Black-body radiationEdit Thermal radiation is a byproduct of the collisions arising from various vibrational motions of atoms. These collisions cause the electrons of the atoms to emit thermal photons (known as black-body radiation). Photons are emitted anytime an electric charge is accelerated (as happens when electron clouds of two atoms collide). Even individual molecules with internal temperatures greater than absolute zero also emit black-body radiation from their atoms. In any bulk quantity of a substance at equilibrium, black-body photons are emitted across a range of wavelengths in a spectrum that has a bell curve-like shape called a Planck curve (see graph in Fig. 5 at right). The top of a Planck curve (the peak emittance wavelength) is located in a particular part of the electromagnetic spectrum depending on the temperature of the black-body. Substances at extreme cryogenic temperatures emit at long radio wavelengths whereas extremely hot temperatures produce short gamma rays (see Table of common temperatures). Black-body radiation diffuses thermal energy throughout a substance as the photons are absorbed by neighboring atoms, transferring momentum in the process. Black-body photons also easily escape from a substance and can be absorbed by the ambient environment; kinetic energy is lost in the process. As established by the Stefan–Boltzmann law, the intensity of black-body radiation increases as the fourth power of absolute temperature. Thus, a black-body at 824 K (just short of glowing dull red) emits 60 times the radiant power as it does at 296 K (room temperature). This is why one can so easily feel the radiant heat from hot objects at a distance. At higher temperatures, such as those found in an incandescent lamp, black-body radiation can be the principal mechanism by which thermal energy escapes a system. Table of thermodynamic temperaturesEdit The full range of the thermodynamic temperature scale, from absolute zero to absolute hot, and some notable points between them are shown in the table below. (precisely by definition) \|0 K\|\|∞ \| \|450 pK\|\|6,400 km\| (precisely by definition) \|0.001 K\|\|2.897 77 m\| (Radio, FM band) \|2.725 K\|\|1.063 mm (peak wavelength)\| \|Water's triple point\|\|273.16 K\|\|10.6083 µm\| (Long wavelength I.R.) \|ISO 1 standard temperature for precision metrology (precisely 20 °C by definition) \|293.15 K\|\|9.88495 µm\| (Long wavelength I.R.) \|Incandescent lamp[A]\|\|2500 K[B]\|\|1.16 µm\| \|Sun’s visible surface[C]\|\|5778 K\|\|501.5 nm\| \|28,000 K\|\|100 nm\| (Far Ultraviolet light) \|Sun's core\|\|16 MK\|\|0.18 nm (X-rays)\| \|350 MK\|\|8.3 × 10−3 nm\| \|Sandia National Labs’ \|2 GK\|\|1.4 × 10−3 nm\| \|Core of a high-mass star on its last day \|3 GK\|\|1 × 10−3 nm\| \|Merging binary neutron star system \|350 GK\|\|8 × 10−6 nm\| \|Gamma-ray burst progenitors\|\|1 TK\|\|3 × 10−6 nm\| \|CERN’s proton vs. \|10 TK\|\|3 × 10−7 nm\| \|Universe 5.391 × 10−44 s after the Big Bang \|1.417 × 1032 K\|\|1.616 × 10−26 nm\| - For a true blackbody (which tungsten filaments are not). Tungsten filaments' emissivity is greater at shorter wavelengths, which makes them appear whiter. - The 2500 K value is approximate. - Effective photosphere temperature. - For a true blackbody (which the plasma was not). The Z machine's dominant emission originated from 40 MK electrons (soft x–ray emissions) within the plasma. The heat of phase changesEdit The kinetic energy of particle motion is just one contributor to the total thermal energy in a substance; another is phase transitions, which are the potential energy of molecular bonds that can form in a substance as it cools (such as during condensing and freezing). The thermal energy required for a phase transition is called latent heat. This phenomenon may more easily be grasped by considering it in the reverse direction: latent heat is the energy required to break chemical bonds (such as during evaporation and melting). Almost everyone is familiar with the effects of phase transitions; for instance, steam at 100 °C can cause severe burns much faster than the 100 °C air from a hair dryer. This occurs because a large amount of latent heat is liberated as steam condenses into liquid water on the skin. Even though thermal energy is liberated or absorbed during phase transitions, pure chemical elements, compounds, and eutectic alloys exhibit no temperature change whatsoever while they undergo them (see Fig. 7, below right). Consider one particular type of phase transition: melting. When a solid is melting, crystal lattice chemical bonds are being broken apart; the substance is transitioning from what is known as a more ordered state to a less ordered state. In Fig. 7, the melting of ice is shown within the lower left box heading from blue to green. At one specific thermodynamic point, the melting point (which is 0 °C across a wide pressure range in the case of water), all the atoms or molecules are, on average, at the maximum energy threshold their chemical bonds can withstand without breaking away from the lattice. Chemical bonds are all-or-nothing forces: they either hold fast, or break; there is no in-between state. Consequently, when a substance is at its melting point, every joule of added thermal energy only breaks the bonds of a specific quantity of its atoms or molecules, converting them into a liquid of precisely the same temperature; no kinetic energy is added to translational motion (which is what gives substances their temperature). The effect is rather like popcorn: at a certain temperature, additional thermal energy can't make the kernels any hotter until the transition (popping) is complete. If the process is reversed (as in the freezing of a liquid), thermal energy must be removed from a substance. As stated above, the thermal energy required for a phase transition is called latent heat. In the specific cases of melting and freezing, it's called enthalpy of fusion or heat of fusion. If the molecular bonds in a crystal lattice are strong, the heat of fusion can be relatively great, typically in the range of 6 to 30 kJ per mole for water and most of the metallic elements. If the substance is one of the monatomic gases, (which have little tendency to form molecular bonds) the heat of fusion is more modest, ranging from 0.021 to 2.3 kJ per mole. Relatively speaking, phase transitions can be truly energetic events. To completely melt ice at 0 °C into water at 0 °C, one must add roughly 80 times the thermal energy as is required to increase the temperature of the same mass of liquid water by one degree Celsius. The metals' ratios are even greater, typically in the range of 400 to 1200 times. And the phase transition of boiling is much more energetic than freezing. For instance, the energy required to completely boil or vaporize water (what is known as enthalpy of vaporization) is roughly 540 times that required for a one-degree increase. Water's sizable enthalpy of vaporization is why one's skin can be burned so quickly as steam condenses on it (heading from red to green in Fig. 7 above). In the opposite direction, this is why one's skin feels cool as liquid water on it evaporates (a process that occurs at a sub-ambient wet-bulb temperature that is dependent on relative humidity). Water's highly energetic enthalpy of vaporization is also an important factor underlying why solar pool covers (floating, insulated blankets that cover swimming pools when not in use) are so effective at reducing heating costs: they prevent evaporation. For instance, the evaporation of just 20 mm of water from a 1.29-meter-deep pool chills its water 8.4 degrees Celsius (15.1 °F). The total energy of all particle motion translational and internal, including that of conduction electrons, plus the potential energy of phase changes, plus zero-point energy comprise the internal energy of a substance. Internal energy at absolute zeroEdit As a substance cools, different forms of internal energy and their related effects simultaneously decrease in magnitude: the latent heat of available phase transitions is liberated as a substance changes from a less ordered state to a more ordered state; the translational motions of atoms and molecules diminish (their kinetic temperature decreases); the internal motions of molecules diminish (their internal temperature decreases); conduction electrons (if the substance is an electrical conductor) travel somewhat slower; and black-body radiation's peak emittance wavelength increases (the photons' energy decreases). When the particles of a substance are as close as possible to complete rest and retain only ZPE-induced quantum mechanical motion, the substance is at the temperature of absolute zero (T = 0). Note that whereas absolute zero is the point of zero thermodynamic temperature and is also the point at which the particle constituents of matter have minimal motion, absolute zero is not necessarily the point at which a substance contains zero internal energy; one must be very precise with what one means by internal energy. Often, all the phase changes that can occur in a substance, will have occurred by the time it reaches absolute zero. However, this is not always the case. Notably, T = 0 helium remains liquid at room pressure (Fig. 9 at right) and must be under a pressure of at least 25 bar (2.5 MPa) to crystallize. This is because helium's heat of fusion (the energy required to melt helium ice) is so low (only 21 joules per mole) that the motion-inducing effect of zero-point energy is sufficient to prevent it from freezing at lower pressures. A further complication is that many solids change their crystal structure to more compact arrangements at extremely high pressures (up to millions of bars, or hundreds of gigapascals). These are known as solid–solid phase transitions wherein latent heat is liberated as a crystal lattice changes to a more thermodynamically favorable, compact one. The above complexities make for rather cumbersome blanket statements regarding the internal energy in T = 0 substances. Regardless of pressure though, what can be said is that at absolute zero, all solids with a lowest-energy crystal lattice such those with a closest-packed arrangement (see Fig. 8, above left) contain minimal internal energy, retaining only that due to the ever-present background of zero-point energy. One can also say that for a given substance at constant pressure, absolute zero is the point of lowest enthalpy (a measure of work potential that takes internal energy, pressure, and volume into consideration). Lastly, it is always true to say that all T = 0 substances contain zero kinetic thermal energy. Practical applications for thermodynamic temperatureEdit Thermodynamic temperature is useful not only for scientists, it can also be useful for lay-people in many disciplines involving gases. By expressing variables in absolute terms and applying Gay–Lussac's law of temperature/pressure proportionality, solutions to everyday problems are straightforward; for instance, calculating how a temperature change affects the pressure inside an automobile tire. If the tire has a cold gage pressure of 200 kPa, then its absolute pressure is 300 kPa. Room temperature ("cold" in tire terms) is 296 K. If the tire temperature is 20 °C hotter (20 kelvins), the solution is calculated as 316 K/296 K = 6.8% greater thermodynamic temperature and absolute pressure; that is, an absolute pressure of 320 kPa, which is a gage pressure of 220 kPa. Relationship to ideal gas lawEdit The thermodynamic temperature is closely linked to the ideal gas law and its consequences. It can be linked also to the second law of thermodynamics. The thermodynamic temperature can be shown to have special properties, and in particular can be seen to be uniquely defined (up to some constant multiplicative factor) by considering the efficiency of idealized heat engines. Thus the ratio T2/T1 of two temperatures T1 and T2 is the same in all absolute scales. Strictly speaking, the temperature of a system is well-defined only if it is at thermal equilibrium. From a microscopic viewpoint, a material is at thermal equilibrium if the quantity of heat between its individual particles cancel out. There are many possible scales of temperature, derived from a variety of observations of physical phenomena. Loosely stated, temperature differences dictate the direction of heat between two systems such that their combined energy is maximally distributed among their lowest possible states. We call this distribution "entropy". To better understand the relationship between temperature and entropy, consider the relationship between heat, work and temperature illustrated in the Carnot heat engine. The engine converts heat into work by directing a temperature gradient between a higher temperature heat source, TH, and a lower temperature heat sink, TC, through a gas filled piston. The work done per cycle is equal in magnitude to net heat taken up, which is sum of the heat qH taken up by the engine from the high-temperature source, plus the waste heat given off by the engine, qC < 0. The efficiency of the engine is the work divided by the heat put into the system or Carnot's theorem states that all reversible engines operating between the same heat reservoirs are equally efficient. Thus, any reversible heat engine operating between temperatures T1 and T2 must have the same efficiency, that is to say, the efficiency is the function of only temperatures In addition, a reversible heat engine operating between temperatures T1 and T3 must have the same efficiency as one consisting of two cycles, one between T1 and another (intermediate) temperature T2, and the second between T2 andT3. If this were not the case, then energy (in the form of q) will be wasted or gained, resulting in different overall efficiencies every time a cycle is split into component cycles; clearly a cycle can be composed of any number of smaller cycles. With this understanding of q1, q2 and q3, mathematically, But since the first function is NOT a function of T2, the product of the final two functions MUST result in the removal of T2 as a variable. The only way is therefore to define the function f as follows: i.e. the ratio of heat exchanged is a function of the respective temperatures at which they occur. We can choose any monotonic function for our ; it is a matter of convenience and convention that we choose . Choosing then one fixed reference temperature (i.e. triple point of water), we establish the thermodynamic temperature scale. Such a definition coincides with that of the ideal gas derivation; also it is this definition of the thermodynamic temperature that enables us to represent the Carnot efficiency in terms of TH and TC, and hence derive that the (complete) Carnot cycle is isentropic: Substituting this back into our first formula for efficiency yields a relationship in terms of temperature: Note that for TC = 0 the efficiency is 100% and that efficiency becomes greater than 100% for TC < 0, which is unrealistic. Subtracting 1 from the right hand side of Equation (4) and the middle portion gives and thus The generalization of this equation is the Clausius theorem, which proposes the existence of a state function (i.e., a function which depends only on the state of the system, not on how it reached that state) defined (up to an additive constant) by where the subscript indicates heat transfer in a reversible process. The function is the entropy of the system, mentioned previously, and the change of around any cycle is zero (as is necessary for any state function). Equation 5 can be rearranged to get an alternative definition for temperature in terms of entropy and heat (to avoid a logic loop, we should first define entropy through statistical mechanics): For a constant-volume system in which the entropy is a function of its energy , and the thermodynamic temperature is therefore given by 1702–1703: Guillaume Amontons (1663–1705) published two papers that may be used to credit him as being the first researcher to deduce the existence of a fundamental (thermodynamic) temperature scale featuring an absolute zero. He made the discovery while endeavoring to improve upon the air thermometers in use at the time. His J-tube thermometers comprised a mercury column that was supported by a fixed mass of air entrapped within the sensing portion of the thermometer. In thermodynamic terms, his thermometers relied upon the volume / temperature relationship of gas under constant pressure. His measurements of the boiling point of water and the melting point of ice showed that regardless of the mass of air trapped inside his thermometers or the weight of mercury the air was supporting, the reduction in air volume at the ice point was always the same ratio. This observation led him to posit that a sufficient reduction in temperature would reduce the air volume to zero. In fact, his calculations projected that absolute zero was equivalent to −240 °C—only 33.15 degrees short of the true value of −273.15 °C. Amonton's discovery of a one-to-one relationship between absolute temperature and absolute pressure was rediscovered a century later and popularized within the scientific community by Joseph Louis Gay-Lussac. Today, this principle of thermodynamics is commonly known as Gay-Lussac's law but is also known as Amonton's law. 1742: Anders Celsius (1701–1744) created a “backwards” version of the modern Celsius temperature scale. In Celsius's original scale, zero represented the boiling point of water and 100 represented the melting point of ice. In his paper Observations of two persistent degrees on a thermometer, he recounted his experiments showing that ice's melting point was effectively unaffected by pressure. He also determined with remarkable precision how water's boiling point varied as a function of atmospheric pressure. He proposed that zero on his temperature scale (water's boiling point) would be calibrated at the mean barometric pressure at mean sea level. 1744: Coincident with the death of Anders Celsius, the famous botanist Carl Linnaeus (1707–1778) effectively reversed Celsius's scale upon receipt of his first thermometer featuring a scale where zero represented the melting point of ice and 100 represented water's boiling point. The custom-made linnaeus-thermometer, for use in his greenhouses, was made by Daniel Ekström, Sweden's leading maker of scientific instruments at the time. For the next 204 years, the scientific and thermometry communities worldwide referred to this scale as the centigrade scale. Temperatures on the centigrade scale were often reported simply as degrees or, when greater specificity was desired, degrees centigrade. The symbol for temperature values on this scale was °C (in several formats over the years). Because the term centigrade was also the French-language name for a unit of angular measurement (one-hundredth of a right angle) and had a similar connotation in other languages, the term "centesimal degree" was used when very precise, unambiguous language was required by international standards bodies such as the International Bureau of Weights and Measures (Bureau international des poids et mesures) (BIPM). The 9th CGPM (General Conference on Weights and Measures (Conférence générale des poids et mesures) and the CIPM (International Committee for Weights and Measures (Comité international des poids et mesures) formally adopted degree Celsius (symbol: °C) in 1948. 1777: In his book Pyrometrie (Berlin: Haude & Spener, 1779) completed four months before his death, Johann Heinrich Lambert (1728–1777), sometimes incorrectly referred to as Joseph Lambert, proposed an absolute temperature scale based on the pressure/temperature relationship of a fixed volume of gas. This is distinct from the volume/temperature relationship of gas under constant pressure that Guillaume Amontons discovered 75 years earlier. Lambert stated that absolute zero was the point where a simple straight-line extrapolation reached zero gas pressure and was equal to −270 °C. Circa 1787: Notwithstanding the work of Guillaume Amontons 85 years earlier, Jacques Alexandre César Charles (1746–1823) is often credited with discovering, but not publishing, that the volume of a gas under constant pressure is proportional to its absolute temperature. The formula he created was V1/T1 = V2/T2. 1802: Joseph Louis Gay-Lussac (1778–1850) published work (acknowledging the unpublished lab notes of Jacques Charles fifteen years earlier) describing how the volume of gas under constant pressure changes linearly with its absolute (thermodynamic) temperature. This behavior is called Charles's Law and is one of the gas laws. His are the first known formulas to use the number 273 for the expansion coefficient of gas relative to the melting point of ice (indicating that absolute zero was equivalent to −273 °C). 1848: William Thomson, (1824–1907) also known as Lord Kelvin, wrote in his paper, On an Absolute Thermometric Scale, of the need for a scale whereby infinite cold (absolute zero) was the scale's zero point, and which used the degree Celsius for its unit increment. Like Gay-Lussac, Thomson calculated that absolute zero was equivalent to −273 °C on the air thermometers of the time. This absolute scale is known today as the kelvin thermodynamic temperature scale. It's noteworthy that Thomson's value of −273 was actually derived from 0.00366, which was the accepted expansion coefficient of gas per degree Celsius relative to the ice point. The inverse of −0.00366 expressed to five significant digits is −273.22 °C which is remarkably close to the true value of −273.15 °C. 1859: Macquorn Rankine (1820–1872) proposed a thermodynamic temperature scale similar to William Thomson's but which used the degree Fahrenheit for its unit increment. This absolute scale is known today as the Rankine thermodynamic temperature scale. 1877–1884: Ludwig Boltzmann (1844–1906) made major contributions to thermodynamics through an understanding of the role that particle kinetics and black body radiation played. His name is now attached to several of the formulas used today in thermodynamics. Circa 1930s: Gas thermometry experiments carefully calibrated to the melting point of ice and boiling point of water showed that absolute zero was equivalent to −273.15 °C. 1948: Resolution 3 of the 9th CGPM (Conférence Générale des Poids et Mesures, also known as the General Conference on Weights and Measures) fixed the triple point of water at precisely 0.01 °C. At this time, the triple point still had no formal definition for its equivalent kelvin value, which the resolution declared "will be fixed at a later date". The implication is that if the value of absolute zero measured in the 1930s was truly −273.15 °C, then the triple point of water (0.01 °C) was equivalent to 273.16 K. Additionally, both the CIPM (Comité international des poids et mesures, also known as the International Committee for Weights and Measures) and the CGPM formally adopted the name Celsius for the degree Celsius and the Celsius temperature scale. 1954: Resolution 3 of the 10th CGPM gave the kelvin scale its modern definition by choosing the triple point of water as its upper defining point (with no change to absolute zero being the null point) and assigning it a temperature of precisely 273.16 kelvins (what was actually written 273.16 degrees Kelvin at the time). This, in combination with Resolution 3 of the 9th CGPM, had the effect of defining absolute zero as being precisely zero kelvins and −273.15 °C. 1967/1968: Resolution 3 of the 13th CGPM renamed the unit increment of thermodynamic temperature kelvin, symbol K, replacing degree absolute, symbol °K. Further, feeling it useful to more explicitly define the magnitude of the unit increment, the 13th CGPM also decided in Resolution 4 that "The kelvin, unit of thermodynamic temperature, is the fraction 1/273.16 of the thermodynamic temperature of the triple point of water". 2005: The CIPM (Comité International des Poids et Mesures, also known as the International Committee for Weights and Measures) affirmed that for the purposes of delineating the temperature of the triple point of water, the definition of the kelvin thermodynamic temperature scale would refer to water having an isotopic composition defined as being precisely equal to the nominal specification of Vienna Standard Mean Ocean Water. 2019: In November 2018, the 26th General Conference on Weights and Measures (CGPM) changed the definition of the Kelvin by fixing the Boltzmann constant to 1.380649×10−23 when expressed in the unit J/K. This change (and other changes in the definition of SI units) was made effective on the 144th anniversary of the Metre Convention, 20 May 2019. - Absolute zero - Hagedorn temperature - Adiabatic process - Boltzmann constant - Carnot heat engine - Energy conversion efficiency - Equipartition theorem - First law of thermodynamics - Gas laws - International System of Quantities - Ideal gas law - Laws of thermodynamics - Maxwell–Boltzmann distribution - Orders of magnitude (temperature) - Phase transition - Planck's law of black-body radiation - Rankine scale - Specific heat capacity - Standard enthalpy change of fusion - Standard enthalpy change of vaporization - Temperature conversion formulas - Thermal radiation - Thermodynamic beta - Thermodynamic equations - Thermodynamic equilibrium - Thermodynamics Category (list of articles) - Timeline of heat engine technology - Timeline of temperature and pressure measurement technology - Triple point - In the following notes, wherever numeric equalities are shown in concise form, such as 1.85487(14)×1043, the two digits between the parentheses denotes the uncertainty at 1-σ (1 standard deviation, 68% confidence level) in the two least significant digits of the significand. - Rankine, W. J. M., "A manual of the steam engine and other prime movers", Richard Griffin and Co., London (1859), p. 306–307. - William Thomson, 1st Baron Kelvin, "Heat", Adam and Charles Black, Edinburgh (1880), p. 39. - absolute zero, they can not fully achieve a state of zero temperature. However, even if scientists could remove all kinetic thermal energy from matter, quantum mechanical zero-point energy (ZPE) causes particle motion that can never be eliminated. Encyclopædia Britannica Online defines zero-point energy as the "vibrational energy that molecules retain even at the absolute zero of temperature". ZPE is the result of all-pervasive energy fields in the vacuum between the fundamental particles of nature; it is responsible for the Casimir effect and other phenomena. See Zero Point Energy and Zero Point Field. See also Solid Helium Archived 2008-02-12 at the Wayback Machine by the University of Alberta's Department of Physics to learn more about ZPE's effect on Bose–Einstein condensates of helium. Although absolute zero (T=0) is not a state of zero molecular motion, it is the point of zero temperature and, in accordance with the Boltzmann constant, is also the point of zero particle kinetic energy and zero kinetic velocity. To understand how atoms can have zero kinetic velocity and simultaneously be vibrating due to ZPE, consider the following thought experiment: two T=0 helium atoms in zero gravity are carefully positioned and observed to have an average separation of 620 pm between them (a gap of ten atomic diameters). It's an "average" separation because ZPE causes them to jostle about their fixed positions. Then one atom is given a kinetic kick of precisely 83 yoctokelvins (1 yK = 1×10−24 K). This is done in a way that directs this atom's velocity vector at the other atom. With 83 yK of kinetic energy between them, the 620 pm gap through their common barycenter would close at a rate of 719 pm/s and they would collide after 0.862 second. This is the same speed as shown in the Fig. 1 animation above. Before being given the kinetic kick, both T=0 atoms had zero kinetic energy and zero kinetic velocity because they could persist indefinitely in that state and relative orientation even though both were being jostled by ZPE. At T=0, no kinetic energy is available for transfer to other systems. The Boltzmann constant and its related formulas describe the realm of particle kinetics and velocity vectors whereas ZPE is an energy field that jostles particles in ways described by the mathematics of quantum mechanics. In atomic and molecular collisions in gases, ZPE introduces a degree of chaos, i.e., unpredictability, to rebound kinetics; it is as likely that there will be less ZPE-induced particle motion after a given collision as more. This random nature of ZPE is why it has no net effect upon either the pressure or volume of any bulk quantity (a statistically significant quantity of particles) of T>0 K gases. However, in T=0 condensed matter; e.g., solids and liquids, ZPE causes inter-atomic jostling where atoms would otherwise be perfectly stationary. Inasmuch as the real-world effects that ZPE has on substances can vary as one alters a thermodynamic system (for example, due to ZPE, helium won't freeze unless under a pressure of at least 25 bar or 2.5 MPa), ZPE is very much a form of thermal energy and may properly be included when tallying a substance's internal energy. Note too that absolute zero serves as the baseline atop which thermodynamics and its equations are founded because they deal with the exchange of thermal energy between "systems" (a plurality of particles and fields modeled as an average). Accordingly, one may examine ZPE-induced particle motion within a system that is at absolute zero but there can never be a net outflow of thermal energy from such a system. Also, the peak emittance wavelength of black-body radiation shifts to infinity at absolute zero; indeed, a peak no longer exists and black-body photons can no longer escape. Because of ZPE, however, virtual photons are still emitted at T=0. Such photons are called "virtual" because they can't be intercepted and observed. Furthermore, this zero-point radiation has a unique zero-point spectrum. However, even though a T=0 system emits zero-point radiation, no net heat flow Q out of such a system can occur because if the surrounding environment is at a temperature greater than T=0, heat will flow inward, and if the surrounding environment is at T=0, there will be an equal flux of ZP radiation both inward and outward. A similar Q equilibrium exists at T=0 with the ZPE-induced spontaneous emission of photons (which is more properly called a stimulated emission in this context). The graph at upper right illustrates the relationship of absolute zero to zero-point energy. The graph also helps in the understanding of how zero-point energy got its name: it is the vibrational energy matter retains at the zero-kelvin point. Derivation of the classical electromagnetic zero-point radiation spectrum via a classical thermodynamic operation involving van der Waals forces, Daniel C. Cole, Physical Review A, 42 (1990) 1847. - CODATA Value: Boltzmann constant. The NIST Reference on Constants, Units, and Uncertainty. National Institute of Standards and Technology. - Georgia State University, HyperPhysics Project, “Equipartition of Energy" - "SI brochure, section 188.8.131.52". International Bureau of Weights and Measures. Archived from the original on 26 September 2007. Retrieved 9 May 2008. - Newell, D B; Cabiati, F; Fischer, J; Fujii, K; Karshenboim, S G; Margolis, H S; de Mirandés, E; Mohr, P J; Nez, F; Pachucki, K; Quinn, T J; Taylor, B N; Wang, M; Wood, B M; Zhang, Z; et al. (Committee on Data for Science and Technology (CODATA) Task Group on Fundamental Constants) (29 January 2018). "The CODATA 2017 values of h, e, k, and NA for the revision of the SI". Metrologia. 55 (1): L13–L16. Bibcode:2018Metro..55L..13N. doi:10.1088/1681-7575/aa950a. - "SI Redefinition – Kelvin: Boltzmann Constant". National Institute of Standards and Technology. Archived from the original on 1 July 2020. Retrieved 13 Dec 2020. - "Accoustic Thermometry". National Institute of Standards and Technology. Archived from the original on 23 September 2020. Retrieved 13 Dec 2020. - At non-relativistic temperatures of less than about 30 GK, classical mechanics are sufficient to calculate the velocity of particles. At 30 GK, individual neutrons (the constituent of neutron stars and one of the few materials in the universe with temperatures in this range) have a 1.0042 γ (gamma or Lorentz factor). Thus, the classic Newtonian formula for kinetic energy is in error less than half a percent for temperatures less than 30 GK. - Even room–temperature air has an average molecular translational speed (not vector-isolated velocity) of 1822 km/hour. This is relatively fast for something the size of a molecule considering there are roughly 2.42×1016 of them crowded into a single cubic millimeter. Assumptions: Average molecular weight of wet air = 28.838 g/mol and T = 296.15 K. Assumption’s primary variables: An altitude of 194 meters above mean sea level (the world–wide median altitude of human habitation), an indoor temperature of 23 °C, a dewpoint of 9 °C (40.85% relative humidity), and 760 mmHg (101.325 kPa) sea level–corrected barometric pressure. - Citation: Adiabatic Cooling of Cesium to 700 nK in an Optical Lattice, A. Kastberg et al., Physical Review Letters 74, No. 9, 27 Feb. 1995, Pg. 1542. It’s noteworthy that a record cold temperature of 450 pK in a Bose–Einstein condensate of sodium atoms (achieved by A. E. Leanhardt et al.. of MIT) equates to an average vector-isolated atom velocity of 0.4 mm/s and an average atom speed of 0.7 mm/s. - The rate of translational motion of atoms and molecules is calculated based on thermodynamic temperature as follows: - is the vector-isolated mean velocity of translational particle motion in m/s - kB (Boltzmann constant) = 1.380649×10−23 J/K - T is the thermodynamic temperature in kelvins - m is the molecular mass of substance in kg/particle - One-trillionth of a kelvin is to one kelvin as two sheets of kitchen aluminum foil (0.04 mm) are to the distance around Earth at the equator. - The internal degrees of freedom of molecules cause their external surfaces to vibrate and can also produce overall spinning motions (what can be likened to the jiggling and spinning of an otherwise stationary water balloon). If one examines a single molecule as it impacts a containers’ wall, some of the kinetic energy borne in the molecule’s internal degrees of freedom can constructively add to its translational motion during the instant of the collision and extra kinetic energy will be transferred into the container’s wall. This would induce an extra, localized, impulse-like contribution to the average pressure on the container. However, since the internal motions of molecules are random, they have an equal probability of destructively interfering with translational motion during a collision with a container’s walls or another molecule. Averaged across any bulk quantity of a gas, the internal thermal motions of molecules have zero net effect upon the temperature, pressure, or volume of a gas. Molecules’ internal degrees of freedom simply provide additional locations where kinetic energy is stored. This is precisely why molecular-based gases have greater specific internal capacity than monatomic gases (where additional internal energy must be added to achieve a given temperature rise). - When measured at constant-volume since different amounts of work must be performed if measured at constant-pressure. Nitrogen’s CvH (100 kPa, 20 °C) equals 20.8 J mol–1 K–1 vs. the monatomic gases, which equal 12.4717 J mol–1 K–1. Citations: W.H. Freeman’s Physical Chemistry, Part 3: Change (422 kB PDF, here), Exercise 21.20b, Pg. 787. Also Georgia State University’s Molar Specific Heats of Gases. - The speed at which thermal energy equalizes throughout the volume of a gas is very rapid. However, since gases have extremely low density relative to solids, the heat flux (the thermal power passing per area) through gases is comparatively low. This is why the dead-air spaces in multi-pane windows have insulating qualities. - Diamond is a notable exception. Highly quantized modes of phonon vibration occur in its rigid crystal lattice. Therefore, not only does diamond have exceptionally poor specific heat capacity, it also has exceptionally high thermal conductivity. - Correlation is 752 (W⋅m−1⋅K−1)/(MS⋅cm), σ = 81, through a 7:1 range in conductivity. Value and standard deviation based on data for Ag, Cu, Au, Al, Ca, Be, Mg, Rh, Ir, Zn, Co, Ni, Os, Fe, Pa, Pt, and Sn. Citation: Data from CRC Handbook of Chemistry and Physics, 1st Student Edition and this link to Web Elements' home page. - The cited emission wavelengths are for true black bodies in equilibrium. In this table, only the sun so qualifies. CODATA recommended value of 2.897771955…×10−3 m⋅K used for Wien displacement law constant b. - A record cold temperature of 450 ±80 pK in a Bose–Einstein condensate (BEC) of sodium (23Na) atoms was achieved in 2003 by researchers at MIT. Citation: Cooling Bose–Einstein Condensates Below 500 Picokelvin, A. E. Leanhardt et al., Science 301, 12 Sept. 2003, Pg. 1515. The thermal velocity of the atoms averaged about 0.4 mm/s. It’s noteworthy that this record’s peak emittance black-body wavelength of 6,400 kilometers is roughly the radius of Earth. - The peak emittance wavelength of 2.897 77 m is a frequency of 103.456 MHz - Measurement was made in 2002 and has an uncertainty of ±3 kelvins. A 1989 measurement produced a value of 5777 ±2.5 K. Citation: Overview of the Sun (Chapter 1 lecture notes on Solar Physics by Division of Theoretical Physics, Dept. of Physical Sciences, University of Helsinki). Download paper (252 kB PDF Archived 2014-08-23 at the Wayback Machine) - The 350 MK value is the maximum peak fusion fuel temperature in a thermonuclear weapon of the Teller–Ulam configuration (commonly known as a “hydrogen bomb”). Peak temperatures in Gadget-style fission bomb cores (commonly known as an “atomic bomb”) are in the range of 50 to 100 MK. Citation: Nuclear Weapons Frequently Asked Questions, 3.2.5 Matter At High Temperatures. Link to relevant Web page. All referenced data was compiled from publicly available sources. - Peak temperature for a bulk quantity of matter was achieved by a pulsed-power machine used in fusion physics experiments. The term “bulk quantity” draws a distinction from collisions in particle accelerators wherein high “temperature” applies only to the debris from two subatomic particles or nuclei at any given instant. The >2 GK temperature was achieved over a period of about ten nanoseconds during “shot Z1137.” In fact, the iron and manganese ions in the plasma averaged 3.58 ±0.41 GK (309 ±35 keV) for 3 ns (ns 112 through 115). Citation: Ion Viscous Heating in a Magnetohydrodynamically Unstable Z Pinch at Over 2 × 109 Kelvin, M. G. Haines et al., Physical Review Letters 96, Issue 7, id. 075003. Link to Sandia’s news release. Archived 2006-07-02 at the Wayback Machine - Core temperature of a high–mass (>8–11 solar masses) star after it leaves the main sequence on the Hertzsprung–Russell diagram and begins the alpha process (which lasts one day) of fusing silicon–28 into heavier elements in the following steps: sulfur–32 → argon–36 → calcium–40 → titanium–44 → chromium–48 → iron–52 → nickel–56. Within minutes of finishing the sequence, the star explodes as a Type II supernova. Citation: Stellar Evolution: The Life and Death of Our Luminous Neighbors (by Arthur Holland and Mark Williams of the University of Michigan). Link to Web site. More informative links can be found here, and here Archived 2011-08-14 at the Wayback Machine, and a concise treatise on stars by NASA is here. Archived July 20, 2015, at the Wayback Machine - Based on a computer model that predicted a peak internal temperature of 30 MeV (350 GK) during the merger of a binary neutron star system (which produces a gamma–ray burst). The neutron stars in the model were 1.2 and 1.6 solar masses respectively, were roughly 20 km in diameter, and were orbiting around their barycenter (common center of mass) at about 390 Hz during the last several milliseconds before they completely merged. The 350 GK portion was a small volume located at the pair’s developing common core and varied from roughly 1 to 7 km across over a time span of around 5 ms. Imagine two city-sized objects of unimaginable density orbiting each other at the same frequency as the G4 musical note (the 28th white key on a piano). It’s also noteworthy that at 350 GK, the average neutron has a vibrational speed of 30% the speed of light and a relativistic mass (m) 5% greater than its rest mass (m0). Citation: Oechslin, R.; Janka, H.- T. (2006). "Torus formation in neutron star mergers and well-localized short gamma-ray bursts". Monthly Notices of the Royal Astronomical Society. 368 (4): 1489–1499. arXiv:astro-ph/0507099v2. Bibcode:2006MNRAS.368.1489O. doi:10.1111/j.1365-2966.2006.10238.x. S2CID 15036056. To view a browser-based summary of the research, click here. - NewScientist: Eight extremes: The hottest thing in the universe, 07 March 2011, which stated “While the details of this process are currently unknown, it must involve a fireball of relativistic particles heated to something in the region of a trillion kelvin[s]” - Citation: How do physicists study particles? Archived 2007-10-11 at the Wayback Machine by CERN. - The Planck frequency equals 1.854 87(14) × 1043 Hz (which is the reciprocal of one Planck time). Photons at the Planck frequency have a wavelength of one Planck length. The Planck temperature of 1.416 79(11) × 1032 K equates to a calculated b /T = λmax wavelength of 2.045 31(16) × 10−26 nm. However, the actual peak emittance wavelength quantizes to the Planck length of 1.616 24(12) × 10−26 nm. - Water's enthalpy of fusion (0 °C, 101.325 kPa) equates to 0.062284 eV per molecule so adding one joule of thermal energy to 0 °C water ice causes 1.0021×1020 water molecules to break away from the crystal lattice and become liquid. - Water's enthalpy of fusion is 6.0095 kJ mol−1 K−1 (0 °C, 101.325 kPa). Citation: Water Structure and Science, Water Properties, Enthalpy of fusion, (0 °C, 101.325 kPa) (by London South Bank University). Link to Web site. The only metals with enthalpies of fusion not in the range of 6–30 J mol−1 K−1 are (on the high side): Ta, W, and Re; and (on the low side) most of the group 1 (alkaline) metals plus Ga, In, Hg, Tl, Pb, and Np. Citation: This link to Web Elements' home page. - Xenon value citation: This link to WebElements' xenon data (available values range from 2.3 to 3.1 kJ/mol). It is also noteworthy that helium's heat of fusion of only 0.021 kJ/mol is so weak of a bonding force that zero-point energy prevents helium from freezing unless it is under a pressure of at least 25 atmospheres. - CRC Handbook of Chemistry and Physics, 1st Student Edition and Web Elements. - H2Ospecific heat capacity, Cp = 0.075327 kJ mol−1 K−1 (25 °C); Enthalpy of fusion = 6.0095 kJ/mol (0 °C, 101.325 kPa); Enthalpy of vaporization (liquid) = 40.657 kJ/mol (100 °C). Citation: Water Structure and Science, Water Properties (by London South Bank University). Link to Web site. - Mobile conduction electrons are delocalized, i.e. not tied to a specific atom, and behave rather like a sort of quantum gas due to the effects of zero-point energy. Consequently, even at absolute zero, conduction electrons still move between atoms at the Fermi velocity of about 1.6×106 m/s. Kinetic thermal energy adds to this speed and also causes delocalized electrons to travel farther away from the nuclei. - No other crystal structure can exceed the 74.048% packing density of a closest-packed arrangement. The two regular crystal lattices found in nature that have this density are hexagonal close packed (HCP) and face-centered cubic (FCC). These regular lattices are at the lowest possible energy state. Diamond is a closest-packed structure with an FCC crystal lattice. Note too that suitable crystalline chemical compounds, although usually composed of atoms of different sizes, can be considered as closest-packed structures when considered at the molecular level. One such compound is the common mineral known as magnesium aluminum spinel (MgAl2O4). It has a face-centered cubic crystal lattice and no change in pressure can produce a lattice with a lower energy state. - Nearly half of the 92 naturally occurring chemical elements that can freeze under a vacuum also have a closest-packed crystal lattice. This set includes beryllium, osmium, neon, and iridium (but excludes helium), and therefore have zero latent heat of phase transitions to contribute to internal energy (symbol: U). In the calculation of enthalpy (formula: H = U + pV), internal energy may exclude different sources of thermal energy (particularly ZPE) depending on the nature of the analysis. Accordingly, all T = 0 closest-packed matter under a perfect vacuum has either minimal or zero enthalpy, depending on the nature of the analysis. Use Of Legendre Transforms In Chemical Thermodynamics, Robert A. Alberty, Pure Appl.Chem., 73 (2001) 1349. - Regarding the spelling "gage" vs. "gauge" in the context of pressures measured relative to atmospheric pressure, the preferred spelling varies by country and even by industry. Further, both spellings are often used within a particular industry or country. Industries in British English-speaking countries typically use the spelling "gauge pressure" to distinguish it from the pressure-measuring instrument, which in the U.K., is spelled pressure gage. For the same reason, many of the largest American manufacturers of pressure transducers and instrumentation use the spelling gage pressure (the convention used here) in their formal documentation to distinguish it from the instrument, which is spelled pressure gauge. (see Honeywell-Sensotec's FAQ page and Fluke Corporation's product search page). - Pressure also must be in absolute terms. The air still in a tire at a gage pressure of 0 kPa expands too as it gets hotter. It's not uncommon for engineers to overlook that one must work in terms of absolute pressure when compensating for temperature. For instance, a dominant manufacturer of aircraft tires published a document on temperature-compensating tire pressure, which used gage pressure in the formula. However, the high gage pressures involved (180 psi; 12.4 bar; 1.24 MPa) means the error would be quite small. With low-pressure automobile tires, where gage pressures are typically around 2 bar (200 kPa), failing to adjust to absolute pressure results in a significant error. Referenced document: Aircraft Tire Ratings (155 kB PDF, here). - A difference of 100 kPa is used here instead of the 101.325 kPa value of one standard atmosphere. In 1982, the International Union of Pure and Applied Chemistry (IUPAC) recommended that for the purposes of specifying the physical properties of substances, the standard pressure (atmospheric pressure) should be defined as precisely 100 kPa (≈ 750.062 Torr). Besides being a round number, this had a very practical effect: relatively few people live and work at precisely sea level; 100 kPa equates to the mean pressure at an altitude of about 112 meters, which is closer to the 194–meter, worldwide median altitude of human habitation. For especially low-pressure or high-accuracy work, true atmospheric pressure must be measured. Citation: IUPAC.org, Gold Book, Standard Pressure - Planck, M. (1945). Treatise on Thermodynamics. Dover Publications. p. §90 & §137. eqs.(39), (40), & (65). - Fermi, E. (1956). Thermodynamics. Dover Publications (still in print). p. 48. - A Brief History of Temperature Measurement and; Uppsala University (Sweden), Linnaeus' thermometer - According to The Oxford English Dictionary (OED), the term "Celsius's thermometer" had been used at least as early as 1797. Further, the term "The Celsius or Centigrade thermometer" was again used in reference to a particular type of thermometer at least as early as 1850. The OED also cites this 1928 reporting of a temperature: "My altitude was about 5,800 metres, the temperature was 28° Celsius". However, dictionaries seek to find the earliest use of a word or term and are not a useful resource as regards the terminology used throughout the history of science. According to several writings of Dr. Terry Quinn CBE FRS, Director of the BIPM (1988–2004), including Temperature Scales from the early days of thermometry to the 21st century (150 kB PDF, here) as well as Temperature (2nd Edition / 1990 / Academic Press / 0125696817), the term Celsius in connection with the centigrade scale was not used whatsoever by the scientific or thermometry communities until after the CIPM and CGPM adopted the term in 1948. The BIPM wasn't even aware that degree Celsius was in sporadic, non-scientific use before that time. It's also noteworthy that the twelve-volume, 1933 edition of OED did not even have a listing for the word Celsius (but did have listings for both centigrade and centesimal in the context of temperature measurement). The 1948 adoption of Celsius accomplished three objectives: - All common temperature scales would have their units named after someone closely associated with them; namely, Kelvin, Celsius, Fahrenheit, Réaumur and Rankine. - Notwithstanding the important contribution of Linnaeus who gave the Celsius scale its modern form, Celsius's name was the obvious choice because it began with the letter C. Thus, the symbol °C that for centuries had been used in association with the name centigrade could continue to be used and would simultaneously inherit an intuitive association with the new name. - The new name eliminated the ambiguity of the term centigrade, freeing it to refer exclusively to the French-language name for the unit of angular measurement.
Enzymes called ____________ form breaks in the DNA molecules toprevent the formation of knots in the DNA helix during replication. Which of the following adds new nucleotides to a growing DNA chain? Why does DNA synthesis only proceed in the 5´ to 3´ direction? Because that is the direction in which the two strands of DNA unzip. The 5´ end of each Okazaki fragment begins with: a separate RNA primer. Primase is the enzyme responsible for: making short strands of RNA at the site of replication initiation. In DNA replication, the lagging strand: is synthesized as a series of Okazaki fragments. Okazaki fragments are joined together by: How are the chromosomes of a eukaryote cell replicated? The linear DNA molecules are replicated from multiple origins of replication bidirectionally. ____________, the ends of eukaryotic chromosomes, shorten with every cell replication event. Cancer cells differ from noncancerous cells in that: All of these. Ribose differs from deoxyribose by having: an extra hydroxyl group. Uracil forms a complementary pair with ____________ in RNA and _____________ in DNA. RNA synthesis is also known as: The codon is found in the: Which of the following serves as an “adapter” in protein synthesis and bridges the gap between mRNA and proteins? During protein synthesis, ribosomes: attach to the mRNA molecule and travel along its length. How is the 4-letter language of nucleic acids converted into the 20-word language of amino acids? The 4 nucleic acid bases combine in 3-letter sequences that define different amino acids. The wobble hypothesis states that: certain tRNA anticodons can pair with more than one codon sequence. A sequence of bases located upstream from a reference point occurs: towards the 5´ end of the mRNA sequence. Initiation of transcription requires: a promoter sequence. Why is only one strand of DNA transcribed into mRNA? Because transcribing both DNA strands would produce different amino acid sequences. Aminoacyl-tRNA synthetases ________ link ________ to their respective tRNA molecules. covalently; amino acids In all organisms, the AUG codon codes for: the initiation of translation. The enzyme peptidyl transferase, which catalyzes the transfer of the polypeptide chain attached to the tRNA in the ____________ site to the aminoacyl-tRNA in the ____________ site, is thought to be a(n) ____________ molecule and not a protein. P; A; rRNA Following peptide bond formation between the amino acid in the A site on the ribosome and the growing polypeptide chain, the tRNA in the A site: forms a peptide bond with A site of the ribosome. A polyribosome is: a complex of many ribosome and an mRNA. Interrupted coding sequences include long sequences of bases that do not code for amino acids. These noncoding sequences, called ____________, are found in ____________ cells. Introns in pre-mRNA are known to: undergo excision, whereby they are spliced out of the message. A gene can now be defined as: a DNA sequence that carries information to produce a specific RNA or protein product. Frameshift mutations result from: the insertion or deletion of one or two base pairs. Genes that encode proteins that are always needed are called: The gene that codes for the repressor protein of the E. coli lactose operon is: The molecular switch that controls gene expression is known as: Lactose induces the transcription of the lactose operon by: binding to the allosteric site of the repressor after being converted to allolactose. Bacterial enzymes that are part of a rarely used catabolic pathway are usually organized into a(n) ________________. A repressible operon is usually controlled by: being turned “on,” usually by the end product of the pathway. In the tryptophan operon, the repressor actively binds to the operator when: tryptophan binds to an allosteric site on the repressor. cAMP levels decrease when ______________________________. This results in ____________ of CAP. glucose levels increase; inhibition In the lac operon, cAMP binds to: catabolite activator protein. The role of CAP in the lac operon is: enhancement of RNA polymerase binding to the promoter. Translational controls regulate: the rate at which ribosomes make proteins. Bacterial gene regulation occurs mainly at the _________________ level. Feedback inhibition is: None of these. Feedback inhibition of the first enzyme of a pathway by the end product of the pathway is an example of: Densely staining regions of highly compacted chromatin that are generally not transcribed are: DNA sequences that are methylated by a cell are usually genes that: A TATA box is seen in ______________ cells and is the site where _______________________. eukaryotic; RNA polymerase binds Upstream promoter elements in eukaryotes are: proteins that inhibit RNA polymerase binding to the promoter. Gene amplification involves: extra replication of genes that specify a certain gene product only in cells needingthis product. Some eukaryotic DNA sequences act as introns in the cells of some tissues and exons in the cells of other tissues. This allows: formation of different types of closely related proteins. Splicing together DNA from 2 different organisms is called: recombinant DNA technology. The modification of the DNA of an organism to produce new genes withnew traits is most properly called: The use of organisms to develop useful products is called: viruses that infect bacteria. A ____________ is required to transfer genes from one organism toanother. Transformation is a process whereby: plasmids are transferred into bacterial cells. DNA ligase links two ____________ DNA fragments by ____________ bonds. The amplification of recombinant plasmids occurs by: the process of growth and division of the host cell. In polymerase chain reaction technology, the two strands of DNA areseparated by: ____________ is a technique that can be used to separate DNA molecules The chain termination method is used in: If a protein-coding gene is identified, its function can be studied by using RNA interference to ______________. turn the gene off Why is insulin produced by genetically engineered E. coli cells superior to insulin obtained from animal sources? It contains human rather than animal sequences, reducing the chances of an allergic response. Introducing a gene for a surface protein produced by a disease-causing agent into a nonpathogenic vector is a method to: produce a recombinant vaccine. The FBI currently determined an individual’s DNA profile using: 13 STR markers. Retroviruses make _________ by the process of ________. DNA copies of RNA; reverse transcription An organism in which foreign genes have been incorporated is called a: Transgenic organisms are used to aid studies of: All of these. _______ is an example of a plant that is genetically modified to produce high quantities of β-carotene for conversion to vitamin A.
The asteroid 16 Psyche (named as such because it was the 16th to be discovered) is believed to be the now-exposed core of a differentiated protoplanet that was smashed apart some billions of years ago. Its composition is generally estimated to be 90 percent metallic and 10 percent silicate rock. It’s thought to be much denser than a typical stony object of equivalent size, and it contains approximately 1 percent of the entire mass of the asteroid belt. Assuming that the core is made of iron and nickel, the total value of the asteroid (if we ignore the impact on market prices) would be ~$10,000 quadrillion dollars. There’s a NASA mission to 16 Psyche expected to launch in 2022 and arrive in 2026, and the Hubble Space Telescope just spent some time surveying the core fragment. Psyche’s surface composition and what it’s made of have implications for the kinds of scientific tests and instruments that would be loaded on the probe we send to study the asteroid. The researchers examined Psyche in the ultraviolet — it’s one of just a handful of asteroids to be examined this way — to see if telltale clues about how the light reflected off its surface could reveal details of the asteroid’s internal composition. The data they found generally supports a high iron composition for the object, but they noted that a relatively small amount of iron mixed into rocky materials like olivine can produce Psyche’s ultraviolet signature. This makes it “difficult to quantify the amount of iron that may present on the surface of Psyche.” Scientists have a number of theories for how Psyche could have formed. It could be the remnant iron core of a protoplanet that was destroyed in a single catastrophic impact. It may be the remnant core of a protoplanet that lost its outer layers in heavy impacts that stripped off the surface but didn’t shatter the entire object. There are a number of theories about its formation, and the research to be done at Psyche will tell us a great deal about the early conditions of the solar system. The asteroid belt doesn’t actually preserve the remnants of the ancient solar system as they once existed; the material present there today has been reworked and reworked in collision after collision. As we send spacecraft to asteroids like Vesta and dwarf planets like Ceres, we’ve begun to reconstruct the patterns of those collisions. We know now, for example, that a massive collision on asteroid Vesta a billion years ago is responsible for the high number of Vesta asteroids found on Earth. These ultraviolet analyses of Psyche don’t confirm that the asteroid is an iron-rich core remnant, but they certainly don’t confirm it isn’t (other data gathered on Psyche, like its gravitational impact on other nearby bodies, suggest a very high density). Psyche is the only known metallic core-like body from any planet currently floating around and available to us to access. Put differently: Psyche may represent a scale model of the Earth’s core as it existed during planetary accretion, before gaining enough mass to become a planet. This process may have failed at Psyche because of higher collision energies. One theory for why the asteroid belt failed to coalesce into a planet is that Jupiter kept gravitational energies too high to support planetary accretion, and the material of the belt was either ejected, pounded to dust, flung inwards (during the Great Heavy Bombardment), or swallowed by Jupiter itself. The goals of the Psyche mission are to determine if Psyche represents a planetary core or a mass of previously unmelted material, to gather information on its age, to examine the compositional differences between the minerals in Psyche’s core and the expected contents of Earth’s, to determine the conditions under which it formed, and to characterize its overall topography. There are currently no serious plans to mine Psyche, but if Earth were to begin to colonize the outer solar system (as shown in TV shows like The Expanse), there’s no question that we’d end up tapping the asteroid for material in one way or another. It’s a rock that could teach us about the earliest days of our solar system, while providing the raw materials we’ll need to create its future. Hubble Finds Exoplanet That Could Mirror Planet Nine The planet, known as HD 106906 b, is 11 times the mass of Jupter, and it orbits the binary stars at a distance of nearly 68 billion miles — 730 times greater than the distance between Earth and the sun. Astronomers believe this frigid world could serve as a proxy to help us understand the hypothetical Planet Nine in our own solar system. Hubble Watches Young Planet Grow With New Imaging Technique This marks the first time scientists have been able to directly observe a still-forming exoplanet. Hubble, Now 31, Snaps Stunning Photo of Volatile Star The photo above shows a star known as AG Carinae, which is one of the brightest in the entire galaxy. Hubble may be getting on in years, but it still impresses. Hubble In Safe Mode Again After Computer Failure [UPDATE] The Hubble Space Telescope has been expanding the bounds of human knowledge for more than thirty years. That's not bad for an orbiting installation built in the 1980s that hasn't gotten a service mission in 12 years. Still, the hardware failures are piling up.
In my capstone class for future secondary math teachers, I ask my students to come up with ideas for engaging their students with different topics in the secondary mathematics curriculum. In other words, the point of the assignment was not to devise a full-blown lesson plan on this topic. Instead, I asked my students to think about three different ways of getting their students interested in the topic in the first place. I plan to share some of the best of these ideas on this blog (after asking my students’ permission, of course). This student submission comes from my former student Alyssa Dalling. Her topic, from Precalculus: the equation of a circle. A. How could you as a teacher create an activity or project that involves your topic? A fun way to engage students and also introduce the standard form of an equation of a circle is the following: - Start by separating the class into groups of 2 or 3 - Pass each group a specific amount of flashcards. (Each group will have the same flashcards) - Each flashcard has a picture of a graphed circle and the equation of that circle in standard form - The students will work together to figure out how the pictures of the circle relate to the equation This will help students understand how different aspects of a circle relate to its standard form equation. The following is an example of a flashcard that could be passed out. C. How has this topic appeared in high culture (art, classical music, theatre, etc.)? Circles have been used through history in many different works of art. One such type is called a tessellation. The word Tessellate means to cover a plane with a pattern in such a way as to leave no region uncovered. So, a tessellation is created when a shape or shapes are repeated over and over again. The pictures above show just a few examples of how circles are used in different types of art. A good way to engage students would be to show them a few examples of tessellations using circles. E. How can technology be used to effectively engage students with this topic? Khan Academy has a really fun resource for using equations to graph circles. At the beginning of class, the teacher could allow students to play around with this program. It allows students to see an equation of a circle in standard form then they would graph the circle. It gives hints as well as the answer when students are ready. The good thing about this is that even if a student goes straight to the answer, they are still trying to identify the connection between the equation of the circle and the answer Khan Academy shows.
What our Subtracting Using Objects & Drawings lesson plan includes Lesson Objectives and Overview: Subtracting Using Objects and Drawings lesson plan helps students understand various vocabulary related to subtraction, as well as the symbols minus sign and equals sign. At the end of the lesson, students will be able to subtract whole numbers using objects and drawings. This lesson is for students in 1st grade and 2nd grade. Every lesson plan provides you with a classroom procedure page that outlines a step-by-step guide to follow. You do not have to follow the guide exactly. The guide helps you organize the lesson and details when to hand out worksheets. It also lists information in the blue box that you might find useful. You will find the lesson objectives, state standards, and number of class sessions the lesson should take to complete in this area. In addition, it describes the supplies you will need as well as what and how you need to prepare beforehand. The supplies you will need for this lesson are 10 clothespins per student, one bowl per group of students, and the lesson handouts. To prepare for this lesson ahead of time, you can copy the materials and gather the supplies. Options for Lesson Included with this lesson is an “Options for Lesson” section that lists a number of suggestions for activities to add to the lesson or substitutions for the ones already in the lesson. One of the suggested adjustments to this lesson is to have students use different starting amounts for the subtraction drop activity. You can also have students use different items at different stations for students to rotate through. Finally, you can use food to illustrate the concepts in the lesson by “taking away” food by eating it and writing math sentences. The teacher notes page includes lines that you can use to add your own notes as you’re preparing for this lesson. SUBTRACTING USING OBJECTS & DRAWINGS LESSON PLAN CONTENT PAGES Subtracting Using Objects and Drawings The Subtracting Using Objects & Drawings lesson plan includes two content pages. The lesson begins by asking students to envision that they have five crackers and eat two. How many crackers do they have left? To solve this problem, you have to use an operation, which is a way to solve a math problem. In math, we have four operations: addition, subtraction, multiplication, and division. In order to find the number of crackers left, we have to use subtraction. When talking about subtraction, we also use the words difference, left, minus, take away, remains, and less. The first step is identifying how many of something you have (in this case, crackers). We started with five crackers, and ate three. If you’ve represented this with a drawing, you can simply cross out three crackers to find out how many are left over. In this case, there are two crackers left. When subtracting, we are taking away a certain amount of something to find the answer. The answer is the amount left over. We also call answers to subtraction problems the difference. When you see a minus sign, that means you have to subtract. An equals sign indicates the answer. To write a subtraction problem as a math sentence, we use numbers. For example, we would write 5 – 3 = 2. The lesson then includes another example that uses socks instead of crackers. There are eight socks, and we lose two in the washing machine. We can write this as 8 – 2 = 6. We can also represent this problem using a drawing, which the lesson shows. SUBTRACTING USING OBJECTS & DRAWINGS LESSON PLAN WORKSHEETS The Subtracting Using Objects & Drawings lesson plan includes four worksheets: an activity worksheet, a practice worksheet, a homework assignment, and a quiz. You can refer to the guide on the classroom procedure page to determine when to hand out each worksheet. SUBTRACTION DROP ACTIVITY WORKSHEET Students will work in groups of two or three for this activity, and you will give each student 10 clothespins (or other objects). Each group of students will stand around a small bowl on the floor, and each student in the group will take turns trying to drop their clothespins into the bowl. They will write subtraction sentences using the results of this activity. For example, if they throw 10 clothespins but only seven go in the bowl, they will write 10 – 7 = 3. They will also mark this on the drawings provided on the worksheet. SUBTRACT PRACTICE WORKSHEET For the practice worksheet, students will solve subtraction problems. The problems include illustrations that students can use for help. SUBTRACTING USING OBJECTS & DRAWINGS HOMEWORK ASSIGNMENT The homework assignment asks the teacher to use 10 – 20 edible items (like grapes, chips, goldfish, or pretzels) to create six subtraction problems for students to solve. For each problem, students should draw a picture and write out a math sentence. This lesson includes a quiz that you can use to test your students’ knowledge and understanding of the lesson material. For the quiz, students will solve subtraction problems. Worksheet Answer Keys This lesson plan includes answer keys for the practice worksheet and the quiz. If you choose to administer the lesson pages to your students via PDF, you will need to save a new file that omits these pages. Otherwise, you can simply print out the applicable pages and keep these as reference for yourself when grading assignments.
Students analyze items from the media to answer mathematical questions related to the article. Exponents and working with large numbers are the underlying mathematical ideas this month. Seán P. Madden and Louis Lim Margaret R. Meyer One of my favorite lessons comes from a problem I first heard posed as an open–ended assessment problem by David Clarke at an NCTM conference years ago: Regarding the reflection “On the Area of a Circle” by Cheng, Tay, and Lee (MT April 2012, vol. 105, no. 8, pp. 564-65), it is possible to prove that one can arrange infinitely many sectors of a circle into a rectangle to show that a circle's area is π2. However, the authors' derivation is invalid because they assume their conclusion by using the area of the circle within their proof. Readers comment on published articles or offer their own ideas. Rina Zazkis, Ilya Sinitsky, and Roza Leikin A familiar relationship—the derivative of the area of a circle equals its circumference—is extended to other shapes and solids. Joshua A. Urich and Elizabeth A. Sasse Students peel oranges to explore the surface area and volume of a sphere. Margaret Rathouz, Christopher Novak, and John Clifford Constructing formulas “from scratch” for calculating geometric measurements of shapes—for example, the area of a triangle—involves reasoning deductively and drawing connections between different methods (Usnick, Lamphere, and Bright 1992). Visual and manipulative models also play a role in helping students understand the underlying mathematics implicit in measurement and make sense of the numbers and operations in formulas. The Grazing Goat problem, familiar to many teachers and students, has several variations. The version presented here provides a rich opportunity for engaging students in a project spanning several weeks. Three solutions are discussed: one suitable for a calculus class, one suitable for a geometry class, and one suitable for a precalculus class. Although we start with a calculus approach, most of the article uses only algebra and geometry concepts. Also discussed are the didactics of using projects to open ever-larger fields of mathematics to students. John M. Livermore The study of the history of mathematics can be an important approach to engaging students in learning mathematics. At the same time, mathematics history provides a context that will help students understand how and why certain types of mathematics were developed and used. The study of p has as rich a history as any mathematical topic.
The world is getting “smarter” every day at the speed of light. Companies are constantly making efforts to use machine learning algorithms to make things easier to fulfill consumer expectations. Some day-to-day examples include end-user devices (through face recognition for unlocking smartphones) or detecting credit card fraud (like triggering alerts for unusual purchases). What makes this world smart and more technology-driven are machine learning (ML) and artificial intelligence (AI). Machine learning is dependent on two types of algorithms – supervised learning and unsupervised learning. One significant difference between the two approaches is one works under the surveillance of labeled data while the other does not require any of it. However, some discrepancies between the two approaches in some particular areas make them different from each other. Supervised learning can be defined as a machine learning approach defined by its use of labeled datasets. These datasets are used to design trained or “supervise” algorithms to classify data or predict outcomes accurately. The models can measure their accuracy and learn over time through labeled inputs and outputs. For example, consider learning the model where the input variable, say X, and the output variable, say Y, is mapped into an algorithm to generate the required results. Inference, Y = f(X). Data mining is performed in supervised learning in two different processes – classification and regression. The classification method is applied to the algorithms that accurately assign test data into specific categories like separating numbers from alphabets. In a real-world example, supervised learning is used to classify spam emails separately from the inbox. The common types of classification algorithms include linear classifiers, support vector machines, decision trees, and random forests. In the regression method, supervised learning uses an algorithm that describes the relationship between dependent and independent variables. The regression method proves helpful in predicting numerical values that are based on a different data point, such as determining sales statistics projections for a given business. Some popular regression algorithms include linear regression, logistic regression, and polynomial regression. The following example will help to understand supervised learning – Consider a basket that is full of different kinds of fresh fruit such as apple, bananas, cherries, grapes. The target here is to arrange similar types of fruits in different baskets based on their individuality. If the machine has already worked on a similar activity, it becomes easy to fulfill the tasks. Based on its previous activity, the machine has already gained the knowledge to perform the task, such as it already knows the shape of each fruit present in the basket; it can easily segregate and arrange the same type of fruits in one basket. Here, machine learning extracts knowledge from previous work obtained in the form of training data in Data Mining terminology. In supervised learning, only the input data is present with no corresponding output variable. At the same time, unsupervised learning uses the hidden or underlying structure for the distribution of the data to learn more details. Unsupervised learning makes use of machine learning algorithms to analyze and cluster unlabeled data sets. These algorithms work on identifying hidden patterns in data without any need for human intervention. Another name for unsupervised learning is knowledge discovery. Unsupervised learning is used to perform three main tasks – clustering, association, and dimensionality reduction. Clustering, a data mining technique, is used to group unlabeled data based on their similarities or differences. Clustering algorithms process raw and unclassified data objects into groups that represent structures or patterns in the information. Clustering algorithms are classified into different types, such as specifically exclusive, overlapping, hierarchical, and probabilistic. Association is another type of unsupervised learning that uses different rules to find a relationship between variables in a given dataset. Such methods are constantly used for market basket analysis and recommendation engines, on the outline of “Customers Who Bought This Item Also Bought” recommendations. - Dimensionality reduction The dimensionality reduction method is used when the number of features or dimensions in a given dataset is very high. It aims to help in reducing the number of data inputs to a limited size while maintaining the data integrity of the datasets as much as possible. Such a method is often used in the pre-processing data stage; for example, autoencoders may use it to remove noise from visual data to improve picture quality. At times, more amount of data usually yields more accurate results. It also has a massive impact on the performance of machine learning algorithms (e.g., overfitting) while It can also make it challenging to visualize datasets. The following example will help to understand unsupervised learning better- Again, consider a basket that is full of different kinds of fresh fruit such as apples, bananas, cherries, and grapes. The target here is to arrange the similar type of fruits in different baskets based on their individuality. In this case, the machine does not have any previous knowledge of fruits. This will be the first time the machine encounters such new objects. This is how the machine then processes the tasks – - Selects any physical characteristic of a particular fruit - Arranges fruits based on their color - Red color: apples and cherry - Green color: bananas and grapes - Now along with color, it will learn about the size too - Red color and big size: apple - Red color and small size: cherry - Green color and big size: bananas - Green color and small size: grapes This is how the task gets completed. Here, no prior information is required, meaning there is no need for any training data. It isn’t easy to choose between both the learnings. The answers to both learnings depend on the situation and the task given to it. This blog provides you brief information on what is supervised and unsupervised learning. Come back to this space in a while to learn more about their difference as well as what could be a midway solution for them. To learn more please visit our latest whitepapers on artificial intelligence here.
This chapter summarises the most important data structures in base R. You’ve probably used many (if not all) of them before, but you may not have thought deeply about how they are interrelated. In this brief overview, I won’t discuss individual types in depth. Instead, I’ll show you how they fit together as a whole. If you need more details, you can find them in R’s documentation. R’s base data structures can be organised by their dimensionality (1d, 2d, or nd) and whether they’re homogeneous (all contents must be of the same type) or heterogeneous (the contents can be of different types). This gives rise to the five data types most often used in data analysis: Almost all other objects are built upon these foundations. In the OO field guide you’ll see how more complicated objects are built of these simple pieces. Note that R has no 0-dimensional, or scalar types. Individual numbers or strings, which you might think would be scalars, are actually vectors of length one. Given an object, the best way to understand what data structures it’s composed of is to use str() is short for structure and it gives a compact, human readable description of any R data structure. Take this short quiz to determine if you need to read this chapter. If the answers quickly come to mind, you can comfortably skip this chapter. You can check your answers in answers. What are the three properties of a vector, other than its contents? What are the four common types of atomic vectors? What are the two rare types? What are attributes? How do you get them and set them? How is a list different from an atomic vector? How is a matrix different from a data frame? Can you have a list that is a matrix? Can a data frame have a column that is a matrix? Vectors introduces you to atomic vectors and lists, R’s 1d data structures. Attributes takes a small detour to discuss attributes, R’s flexible metadata specification. Here you’ll learn about factors, an important data structure created by setting attributes of an atomic vector. Matrices and arrays introduces matrices and arrays, data structures for storing 2d and higher dimensional data. Data frames teaches you about the data frame, the most important data structure for storing data in R. Data frames combine the behaviour of lists and matrices to make a structure ideally suited for the needs of statistical data. The basic data structure in R is the vector. Vectors come in two flavours: atomic vectors and lists. They have three common properties: typeof(), what it is. length(), how many elements it contains. attributes(), additional arbitrary metadata. They differ in the types of their elements: all elements of an atomic vector must be the same type, whereas the elements of a list can have different types. is.vector() does not test if an object is a vector. Instead it returns TRUE only if the object is a vector with no attributes apart from names. Use is.atomic(x) \|\| is.list(x) to test if an object is actually a vector. There are four common types of atomic vectors that I’ll discuss in detail: logical, integer, double (often called numeric), and character. There are two rare types that I will not discuss further: complex and raw. Atomic vectors are usually created with c(), short for combine: dbl_var <- c(1, 2.5, 4.5) # With the L suffix, you get an integer rather than a double int_var <- c(1L, 6L, 10L) # Use TRUE and FALSE (or T and F) to create logical vectors log_var <- c(TRUE, FALSE, T, F) chr_var <- c("these are", "some strings") Atomic vectors are always flat, even if you nest c(1, c(2, c(3, 4))) #> 1 2 3 4 # the same as c(1, 2, 3, 4) #> 1 2 3 4 Missing values are specified with NA, which is a logical vector of length 1. NA will always be coerced to the correct type if used inside c(), or you can create NAs of a specific type with NA_real_ (a double vector), Types and tests Given a vector, you can determine its type with typeof(), or check if it’s a specific type with an “is” function: is.logical(), or, more generally, int_var <- c(1L, 6L, 10L) typeof(int_var) #> "integer" is.integer(int_var) #> TRUE is.atomic(int_var) #> TRUE dbl_var <- c(1, 2.5, 4.5) typeof(dbl_var) #> "double" is.double(dbl_var) #> TRUE is.atomic(dbl_var) #> TRUE is.numeric() is a general test for the “numberliness” of a vector and returns TRUE for both integer and double vectors. It is not a specific test for double vectors, which are often called numeric. is.numeric(int_var) #> TRUE is.numeric(dbl_var) #> TRUE All elements of an atomic vector must be the same type, so when you attempt to combine different types they will be coerced to the most flexible type. Types from least to most flexible are: logical, integer, double, and character. For example, combining a character and an integer yields a character: str(c("a", 1)) #> chr [1:2] "a" "1" When a logical vector is coerced to an integer or double, TRUE becomes 1 and FALSE becomes 0. This is very useful in conjunction with x <- c(FALSE, FALSE, TRUE) as.numeric(x) #> 0 0 1 # Total number of TRUEs sum(x) #> 1 # Proportion that are TRUE mean(x) #> 0.3333333 Coercion often happens automatically. Most mathematical functions ( abs, etc.) will coerce to a double or integer, and most logical operations ( any, etc) will coerce to a logical. You will usually get a warning message if the coercion might lose information. If confusion is likely, explicitly coerce with Lists are different from atomic vectors because their elements can be of any type, including lists. You construct lists by using list() instead of x <- list(1:3, "a", c(TRUE, FALSE, TRUE), c(2.3, 5.9)) str(x) #> List of 4 #> $ : int [1:3] 1 2 3 #> $ : chr "a" #> $ : logi [1:3] TRUE FALSE TRUE #> $ : num [1:2] 2.3 5.9 Lists are sometimes called recursive vectors, because a list can contain other lists. This makes them fundamentally different from atomic vectors. x <- list(list(list(list()))) str(x) #> List of 1 #> $ :List of 1 #> ..$ :List of 1 #> .. ..$ : list() is.recursive(x) #> TRUE c() will combine several lists into one. If given a combination of atomic vectors and lists, c() will coerce the vectors to lists before combining them. Compare the results of x <- list(list(1, 2), c(3, 4)) y <- c(list(1, 2), c(3, 4)) str(x) #> List of 2 #> $ :List of 2 #> ..$ : num 1 #> ..$ : num 2 #> $ : num [1:2] 3 4 str(y) #> List of 4 #> $ : num 1 #> $ : num 2 #> $ : num 3 #> $ : num 4 typeof() a list is list. You can test for a list with is.list() and coerce to a list with as.list(). You can turn a list into an atomic vector with unlist(). If the elements of a list have different types, unlist() uses the same coercion rules as Lists are used to build up many of the more complicated data structures in R. For example, both data frames (described in data frames) and linear models objects (as produced by lm()) are lists: is.list(mtcars) #> TRUE mod <- lm(mpg ~ wt, data = mtcars) is.list(mod) #> TRUE What are the six types of atomic vector? How does a list differ from an atomic vector? is.numeric()fundamentally different to Test your knowledge of vector coercion rules by predicting the output of the following uses of c(1, FALSE) c("a", 1) c(list(1), "a") c(TRUE, 1L) Why do you need to use unlist()to convert a list to an atomic vector? Why doesn’t 1 == "1"true? Why is -1 < FALSEtrue? Why is "one" < 2false? Why is the default missing value, NA, a logical vector? What’s special about logical vectors? (Hint: think about All objects can have arbitrary additional attributes, used to store metadata about the object. Attributes can be thought of as a named list (with unique names). Attributes can be accessed individually with attr() or all at once (as a list) with y <- 1:10 attr(y, "my_attribute") <- "This is a vector" attr(y, "my_attribute") #> "This is a vector" str(attributes(y)) #> List of 1 #> $ my_attribute: chr "This is a vector" structure() function returns a new object with modified attributes: structure(1:10, my_attribute = "This is a vector") #> 1 2 3 4 5 6 7 8 9 10 #> attr(,"my_attribute") #> "This is a vector" By default, most attributes are lost when modifying a vector: attributes(y) #> NULL attributes(sum(y)) #> NULL The only attributes not lost are the three most important: Names, a character vector giving each element a name, described in names. Dimensions, used to turn vectors into matrices and arrays, described in matrices and arrays. Class, used to implement the S3 object system, described in S3. Each of these attributes has a specific accessor function to get and set values. When working with these attributes, use attr(x, "dim"), and You can name a vector in three ways: When creating it: x <- c(a = 1, b = 2, c = 3). By modifying an existing vector in place: x <- 1:3; names(x) <- c("a", "b", "c"). By creating a modified copy of a vector: x <- setNames(1:3, c("a", "b", "c")). Names don’t have to be unique. However, character subsetting, described in subsetting, is the most important reason to use names and it is most useful when the names are unique. Not all elements of a vector need to have a name. If some names are missing, names() will return an empty string for those elements. If all names are missing, names() will return y <- c(a = 1, 2, 3) names(y) #> "a" "" "" z <- c(1, 2, 3) names(z) #> NULL You can create a new vector without names using unname(x), or remove names in place with names(x) <- NULL. One important use of attributes is to define factors. A factor is a vector that can contain only predefined values, and is used to store categorical data. Factors are built on top of integer vectors using two attributes: the class(), “factor”, which makes them behave differently from regular integer vectors, and the levels(), which defines the set of allowed values. x <- factor(c("a", "b", "b", "a")) x #> a b b a #> Levels: a b class(x) #> "factor" levels(x) #> "a" "b" # You can't use values that are not in the levels x <- "c" #> Warning in `[<-.factor`(`tmp`, 2, value = "c"): invalid factor level, NA #> generated x #> a <NA> b a #> Levels: a b # NB: you can't combine factors c(factor("a"), factor("b")) #> 1 1 Factors are useful when you know the possible values a variable may take, even if you don’t see all values in a given dataset. Using a factor instead of a character vector makes it obvious when some groups contain no observations: sex_char <- c("m", "m", "m") sex_factor <- factor(sex_char, levels = c("m", "f")) table(sex_char) #> sex_char #> m #> 3 table(sex_factor) #> sex_factor #> m f #> 3 0 Sometimes when a data frame is read directly from a file, a column you’d thought would produce a numeric vector instead produces a factor. This is caused by a non-numeric value in the column, often a missing value encoded in a special way like -. To remedy the situation, coerce the vector from a factor to a character vector, and then from a character to a double vector. (Be sure to check for missing values after this process.) Of course, a much better plan is to discover what caused the problem in the first place and fix that; using the na.strings argument to read.csv() is often a good place to start. # Reading in "text" instead of from a file here: z <- read.csv(text = "value\n12\n1\n.\n9") typeof(z$value) #> "integer" as.double(z$value) #> 3 2 1 4 # Oops, that's not right: 3 2 1 4 are the levels of a factor, # not the values we read in! class(z$value) #> "factor" # We can fix it now: as.double(as.character(z$value)) #> Warning: NAs introduced by coercion #> 12 1 NA 9 # Or change how we read it in: z <- read.csv(text = "value\n12\n1\n.\n9", na.strings=".") typeof(z$value) #> "integer" class(z$value) #> "integer" z$value #> 12 1 NA 9 # Perfect! :) Unfortunately, most data loading functions in R automatically convert character vectors to factors. This is suboptimal, because there’s no way for those functions to know the set of all possible levels or their optimal order. Instead, use the argument stringsAsFactors = FALSE to suppress this behaviour, and then manually convert character vectors to factors using your knowledge of the data. A global option, options(stringsAsFactors = FALSE), is available to control this behaviour, but I don’t recommend using it. Changing a global option may have unexpected consequences when combined with other code (either from packages, or code that you’re source()ing), and global options make code harder to understand because they increase the number of lines you need to read to understand how a single line of code will behave. While factors look (and often behave) like character vectors, they are actually integers. Be careful when treating them like strings. Some string methods (like grepl()) will coerce factors to strings, while others (like nchar()) will throw an error, and still others (like c()) will use the underlying integer values. For this reason, it’s usually best to explicitly convert factors to character vectors if you need string-like behaviour. In early versions of R, there was a memory advantage to using factors instead of character vectors, but this is no longer the case. An early draft used this code to illustrate structure(1:5, comment = "my attribute") #> 1 2 3 4 5 But when you print that object you don’t see the comment attribute. Why? Is the attribute missing, or is there something else special about it? (Hint: try using help.) What happens to a factor when you modify its levels? f1 <- factor(letters) levels(f1) <- rev(levels(f1)) What does this code do? How do f2 <- rev(factor(letters)) f3 <- factor(letters, levels = rev(letters)) Matrices and arrays dim() attribute to an atomic vector allows it to behave like a multi-dimensional array. A special case of the array is the matrix, which has two dimensions. Matrices are used commonly as part of the mathematical machinery of statistics. Arrays are much rarer, but worth being aware of. Matrices and arrays are created with array(), or by using the assignment form of # Two scalar arguments to specify rows and columns a <- matrix(1:6, ncol = 3, nrow = 2) # One vector argument to describe all dimensions b <- array(1:12, c(2, 3, 2)) # You can also modify an object in place by setting dim() c <- 1:6 dim(c) <- c(3, 2) c #> [,1] [,2] #> [1,] 1 4 #> [2,] 2 5 #> [3,] 3 6 dim(c) <- c(2, 3) c #> [,1] [,2] [,3] #> [1,] 1 3 5 #> [2,] 2 4 6 names() have high-dimensional generalisations: ncol()for matrices, and colnames()for matrices, and dimnames(), a list of character vectors, for arrays. length(a) #> 6 nrow(a) #> 2 ncol(a) #> 3 rownames(a) <- c("A", "B") colnames(a) <- c("a", "b", "c") a #> a b c #> A 1 3 5 #> B 2 4 6 length(b) #> 12 dim(b) #> 2 3 2 dimnames(b) <- list(c("one", "two"), c("a", "b", "c"), c("A", "B")) b #> , , A #> #> a b c #> one 1 3 5 #> two 2 4 6 #> #> , , B #> #> a b c #> one 7 9 11 #> two 8 10 12 c() generalises to rbind() for matrices, and to abind() (provided by the abind package) for arrays. You can transpose a matrix with t(); the generalised equivalent for arrays is You can test if an object is a matrix or array using is.array(), or by looking at the length of the as.array() make it easy to turn an existing vector into a matrix or array. Vectors are not the only 1-dimensional data structure. You can have matrices with a single row or single column, or arrays with a single dimension. They may print similarly, but will behave differently. The differences aren’t too important, but it’s useful to know they exist in case you get strange output from a function ( tapply() is a frequent offender). As always, use str() to reveal the differences. str(1:3) # 1d vector #> int [1:3] 1 2 3 str(matrix(1:3, ncol = 1)) # column vector #> int [1:3, 1] 1 2 3 str(matrix(1:3, nrow = 1)) # row vector #> int [1, 1:3] 1 2 3 str(array(1:3, 3)) # "array" vector #> int [1:3(1d)] 1 2 3 While atomic vectors are most commonly turned into matrices, the dimension attribute can also be set on lists to make list-matrices or list-arrays: l <- list(1:3, "a", TRUE, 1.0) dim(l) <- c(2, 2) l #> [,1] [,2] #> [1,] Integer,3 TRUE #> [2,] "a" 1 These are relatively esoteric data structures, but can be useful if you want to arrange objects into a grid-like structure. For example, if you’re running models on a spatio-temporal grid, it might be natural to preserve the grid structure by storing the models in a 3d array. dim()return when applied to a vector? TRUE, what will How would you describe the following three objects? What makes them different to x1 <- array(1:5, c(1, 1, 5)) x2 <- array(1:5, c(1, 5, 1)) x3 <- array(1:5, c(5, 1, 1)) A data frame is the most common way of storing data in R, and if used systematically makes data analysis easier. Under the hood, a data frame is a list of equal-length vectors. This makes it a 2-dimensional structure, so it shares properties of both the matrix and the list. This means that a data frame has colnames() are the same thing. The length() of a data frame is the length of the underlying list and so is the same as nrow() gives the number of rows. As described in subsetting, you can subset a data frame like a 1d structure (where it behaves like a list), or a 2d structure (where it behaves like a matrix). You create a data frame using data.frame(), which takes named vectors as input: df <- data.frame(x = 1:3, y = c("a", "b", "c")) str(df) #> 'data.frame': 3 obs. of 2 variables: #> $ x: int 1 2 3 #> $ y: Factor w/ 3 levels "a","b","c": 1 2 3 data.frame()’s default behaviour which turns strings into factors. Use stringAsFactors = FALSE to suppress this behaviour: df <- data.frame( x = 1:3, y = c("a", "b", "c"), stringsAsFactors = FALSE) str(df) #> 'data.frame': 3 obs. of 2 variables: #> $ x: int 1 2 3 #> $ y: chr "a" "b" "c" Testing and coercion data.frame is an S3 class, its type reflects the underlying vector used to build it: the list. To check if an object is a data frame, use class() or test explicitly with typeof(df) #> "list" class(df) #> "data.frame" is.data.frame(df) #> TRUE You can coerce an object to a data frame with A vector will create a one-column data frame. A list will create one column for each element; it’s an error if they’re not all the same length. A matrix will create a data frame with the same number of columns and rows as the matrix. Combining data frames You can combine data frames using cbind(df, data.frame(z = 3:1)) #> x y z #> 1 1 a 3 #> 2 2 b 2 #> 3 3 c 1 rbind(df, data.frame(x = 10, y = "z")) #> x y #> 1 1 a #> 2 2 b #> 3 3 c #> 4 10 z When combining column-wise, the number of rows must match, but row names are ignored. When combining row-wise, both the number and names of columns must match. Use plyr::rbind.fill() to combine data frames that don’t have the same columns. It’s a common mistake to try and create a data frame by cbind()ing vectors together. This doesn’t work because cbind() will create a matrix unless one of the arguments is already a data frame. Instead use bad <- data.frame(cbind(a = 1:2, b = c("a", "b"))) str(bad) #> 'data.frame': 2 obs. of 2 variables: #> $ a: Factor w/ 2 levels "1","2": 1 2 #> $ b: Factor w/ 2 levels "a","b": 1 2 good <- data.frame(a = 1:2, b = c("a", "b"), stringsAsFactors = FALSE) str(good) #> 'data.frame': 2 obs. of 2 variables: #> $ a: int 1 2 #> $ b: chr "a" "b" The conversion rules for cbind() are complicated and best avoided by ensuring all inputs are of the same type. Since a data frame is a list of vectors, it is possible for a data frame to have a column that is a list: df <- data.frame(x = 1:3) df$y <- list(1:2, 1:3, 1:4) df #> x y #> 1 1 1, 2 #> 2 2 1, 2, 3 #> 3 3 1, 2, 3, 4 However, when a list is given to data.frame(), it tries to put each item of the list into its own column, so this fails: data.frame(x = 1:3, y = list(1:2, 1:3, 1:4)) #> Error in data.frame(1:2, 1:3, 1:4, check.names = FALSE, stringsAsFactors = TRUE): arguments imply differing number of rows: 2, 3, 4 A workaround is to use I(), which causes data.frame() to treat the list as one unit: dfl <- data.frame(x = 1:3, y = I(list(1:2, 1:3, 1:4))) str(dfl) #> 'data.frame': 3 obs. of 2 variables: #> $ x: int 1 2 3 #> $ y:List of 3 #> ..$ : int 1 2 #> ..$ : int 1 2 3 #> ..$ : int 1 2 3 4 #> ..- attr(*, "class")= chr "AsIs" dfl[2, "y"] #> [] #> 1 2 3 I() adds the AsIs class to its input, but this can usually be safely ignored. Similarly, it’s also possible to have a column of a data frame that’s a matrix or array, as long as the number of rows matches the data frame: dfm <- data.frame(x = 1:3, y = I(matrix(1:9, nrow = 3))) str(dfm) #> 'data.frame': 3 obs. of 2 variables: #> $ x: int 1 2 3 #> $ y: 'AsIs' int [1:3, 1:3] 1 2 3 4 5 6 7 8 9 dfm[2, "y"] #> [,1] [,2] [,3] #> [1,] 2 5 8 Use list and array columns with caution: many functions that work with data frames assume that all columns are atomic vectors. What attributes does a data frame possess? as.matrix()do when applied to a data frame with columns of different types? Can you have a data frame with 0 rows? What about 0 columns? The three properties of a vector are type, length, and attributes. The four common types of atomic vector are logical, integer, double (sometimes called numeric), and character. The two rarer types are complex and raw. Attributes allow you to associate arbitrary additional metadata to any object. You can get and set individual attributes with attr(x, "y") <- value; or get and set all attributes at once with The elements of a list can be any type (even a list); the elements of an atomic vector are all of the same type. Similarly, every element of a matrix must be the same type; in a data frame, the different columns can have different types. You can make “list-array” by assuming dimensions to a list. You can make a matrix a column of a data frame with df$x <- matrix(), or using I()when creating a new data frame data.frame(x = I(matrix())).
Huygen’s wave theory of light cannot explain (a) photoelectric effect If a source of light is moved away from a stationary observer, then the frequency of light wave appears to change because of: (c) doppler’s effect (d) all of these When seen in green light, the saffron and green portions of our National Flag will appear (c) black and green respectively. (d) green and yellow respectively. Two sources of light are said to be coherent when both give out light waves of the same: (a) amplitude and phase (b) intensity and wavelength (d) wavelength and a constant phase difference The locus of all points which oscillates in phase is called __________. (d) none of these The penetration of light into the region of geometrical shadow is called _________. Two waves have intensity ratio 25:4. What is the ratio of maximum to minimum intensity? The fringe width (β) of a diffraction pattern and the slit width d are related as: (a) β ∝ d (b) β ∝1/d (c) β ∝ √d (d) β ∝ 1/d2 Define Red shift and Blue shift in Doppler’s effect. When monochromatic light is incident on a surface separating two media, the reflected and refracted light both have the same frequency as the incident frequency. Explain why? Two slits are made one millimetre apart and the screen is placed one metre away. What is the fringe separation when blue green light of wavelength 500 nm is used? What is the value if refractive index of a medium of polarizing angle 600? For what distance is ray optics a good approximation when the aperture is 3 mm wide and the wavelength is 500 nm? A slit of width 3 mm is illuminated by light of λ = 600 nm at normal incidence. If the distance of the screen from the slit is 60 cm, calculate the distance between the first order minimum on both sides of central maximum. Why is interference pattern nit detected, when two coherent sources are far apart? In Young’s experiment, the width of the fringes obtained with light of wavelength 6000Ǻ is 2 mm. Calculate the fringe width if the entire apparatus is immersed in a liquid medium of refractive index 1.33. What are coherent sources? How does the width of interference fringes in young’s double-slit experiment change when? (a) The distance between the slits and screen is decreased? (b) Frequency of the source is increased? Justify your answer in each case. The maximum intensity in Young’s double slit experiment is I0. Distance between the slits is d= 5λ, where λ is the wavelength of monochromatic light used in the experiment. What will be the intensity of light in front of one of the slits on a screen at a distance D= 10d? A double slit is illuminated by light of wavelength 6000Ǻ. The slits are 0.1 cm apart and the screen is placed 1 m away. Calculate (i) angular position of 10th maximum in radian (ii) separation of two adjacent minima. In a young’s experiment, the width of the fringes obtained with light of wavelength 6000Ǻ is 2 mm. What will be the fringe width, if the entire apparatus is immersed in a liquid of refractive index 1.33? A slit of width ‘d’ is illuminated by white light. For what value of ‘d’ is the first minimum for red light of λ= 650 nm located at point P. For what value of the wavelength of light will be the first diffraction maxima also fall at P?
Location theory, in economics and geography, theory concerned with the geographic location of economic activity; it has become an integral part of economic geography, regional science, and spatial economics. Location theory addresses the questions of what economic activities are located where and why. The location of economic activities can be determined on a broad level such as a region or metropolitan area, or on a narrow one such as a zone, neighbourhood, city block, or an individual site. Johann Heinrich von Thünen, a Prussian landowner, introduced an early theory of agricultural location in Der isolierte Staat (1826) (The Isolated State). The Thünen model suggests that accessibility to the market (town) can create a complete system of agricultural land use. His model envisaged a single market surrounded by farmland, both situated on a plain of complete physical homogeneity. Transportation costs over the plain are related only to the distance traveled and the volume shipped. The model assumes that farmers surrounding the market will produce crops which have the highest market value (highest rent) that will give them the maximum net profit (the location, or land, rent). The determining factor in the location rent will be the transportation costs. When transportation costs are low, the location rent will be high, and vice versa. This situation produces a rent gradient along which the location rent decreases with distance from the market, eventually reaching zero. The Thünen model also addressed the location of intensive versus extensive agriculture in relation to the same market. Intensive agriculture will possess a steep gradient and will locate closer to the market than extensive agriculture. Different crops will possess different rent gradients. Perishable crops (vegetables and dairy products) will possess steep gradients while less perishable crops (grains) will possess less steep gradients. In 1909 the German location economist Alfred Weber formulated a theory of industrial location in his book entitled Über den Standort der Industrien (Theory of the Location of Industries, 1929). Weber’s theory, called the location triangle, sought the optimum location for the production of a good based on the fixed locations of the market and two raw material sources, which geographically form a triangle. He sought to determine the least-cost production location within the triangle by figuring the total costs of transporting raw material from both sites to the production site and product from the production site to the market. The weight of the raw materials and the final commodity are important determinants of the transport costs and the location of production. Commodities that lose mass during production can be transported less expensively from the production site to the market than from the raw material site to the production site. The production site, therefore, will be located near the raw material sources. Where there is no great loss of mass during production, total transportation costs will be lower when located near the market. Once a least-transport-cost location had been established within the triangle, Weber attempted to determine a cheap-labour alternate location. First he plotted the variation of transportation costs against the least-transport-cost location. Next he identified sites around the triangle that had lower labour costs than did the least-transport-cost location. If the transport costs were lower than the labour costs, then a cheap-labour alternative location was determined. Another major contribution to location theory was Walter Christaller’s formulation of the central place theory, which offered geometric explanations as to how settlements and places are located in relation to one another and why settlements function as hamlets, villages, towns, or cities. Test Your Knowledge William Alonso (Location and Land Use: Toward a General Theory of Land Rent, 1964) built upon the Thünen model to account for intra-urban variations in land use. He attempted to apply accessibility requirements to the city centre for various types of land use (housing, commercial, and industry). According to his theory, each land use type has its own rent gradient or bid rent curve. The curve sets the maximum amount of rent any land use type will yield for a specific location. Households, commercial establishments, and industries compete for locations according to each individual bid rent curve and their requirements for access to the city centre. All households will attempt to occupy as much land as possible while staying within their accessibility requirements. Since land is cheaper at the fringe of the city, households with less need for city centre accessibility will locate near the fringe; these will usually be wealthy households. Poor households require greater accessibility to the city centre and therefore will locate near the centre, competing with commercial and industrial establishments. This will tend to create a segregated land use system, because households will not pay commercial and industrial land prices for central locations. The Thünen, Weber, Alonso, and Christaller models are not the sole contributors to location theory, but they are its foundation. These theories have been expanded upon and refined by geographers, economists, and regional scientists.
A statement or an expression in any programming language could involve some arithmetic calculations with more than one operation. For instance, consider the following statement: It has two operators (+ and * respectively), and in order to evaluate the final result, you need to perform these two operations. You can either perform the addition operation first, followed by the multiplication operation. Or you can do it the other way round by performing the multiplication first, followed by the addition operation. In both cases, you will get different results due to the order in which each operation is executed. So how do you determine the correct sequence of operations to produce consistent results in any arithmetic calculation? Precedence Level for Common Operators Note that precedence level varies from 1 to 19. The general rule is that an operator that has a lower precedence level is given higher priority and is said to have higher precedence than the one having a higher precedence level. You can interpret precedence level as rank in priority level of each operator. \|Division\|\|/\|\|3\|\|Left to Right\| \|Multiplication\|\|\|\|3\|\|Left to Right\| \|Modulus\|\|%\|\|3\|\|Left to Right\| \|Addition\|\|+\|\|4\|\|Left to Right\| \|Subtraction\|\|-\|\|4\|\|Left to Right\| \|Assignment\|\|=\|\|14\|\|Right to left\| \|Exponentiation\|\|\|\|14\|\|Right to left\| So if we look at the above table, multiplication will have higher precedence than addition. Similarly, the division will have higher precedence than subtraction. However, addition and subtraction have the same precedence. Similarly, multiplication and division also have the same precedence. We also have a column for associativity. So let's understand what it means. Associativity defines the direction to be followed when executing the operations in a statement. It comes into the picture when statements contain operators that have the same precedence. In order to understand completely why associativity is even needed in the first place, let's take a look at the following example. Need for Associativity Consider the below expression. We know both multiplication and division have the same precedence. So how should you evaluate the expression from the left side, i.e., do you perform division first? In that case, as you can see, the output would be 4. But if you evaluate the expression from the right side, ie, perform multiplication first and then division: The output would be 1. Since we get different outputs in both cases, we need the rule to resolve this ambiguity. That's exactly where associativity comes in. Left to Right Associativity The associativity of an operator describes the direction in which the operations would get executed within a statement. Right to Left Associativity In contrast to the left to right associativity, some operators may have right to left associativity. In the table we saw earlier, the assignment operation had the right to left associativity. Let's consider the following example where a value is assigned to three variables. Since the assignment follows right to left associativity, 25 is first assigned to the variable z. Then, the value of variable z (i.e., now 25) is assigned to variable y. Finally, the value of variable y, which is now 25, is assigned to variable x. Here's how the right to left associativity will break down the above statement: It's simple to look at the above statement and conclude that we're only assigning the same value to three variables x, y, and z. However, it's imperative to understand how the assignment is actually happening. If you think the result would be the same even if we evaluate the expression from left to right, think again! Let's say we evaluate the expression from the left side first. Here's how we'll break the statement down: First, the value of y will be assigned to x. But since y doesn't have a value yet, x will get undefined. Similarly, y will also get undefined in the next step since z also doesn't have any value yet. Finally, z will be assigned a value 25. But at the end, both x and y will have the value undefined in contrast to the right to left evaluation, where both had a value 25. Let's say you're evaluating the following expression. The foremost thing to do is to identify which operators have the highest precedence. We already know multiplication and division have higher precedence than addition and subtraction. So now, we'll evaluate the division and multiplication operations first. Since their associativity tells us that the expression must be evaluated from left to right, we'll evaluate division first. Followed by multiplication: Once we have narrowed down the expression, we'll again look at the precedence of the operators to determine which operation needs to execute first. And then, if there's a conflict in the precedence, we let their associativity break the tie. We already know the precedence and associativity of common operators, but what about the other remaining operators? Here's a table that lays down each operator, along with it's precedence value and associativity type: \|Operator\|\|Operator Use\|\|Operator Associativity\|\|Operator Precedence\| \|0\|\|Method/function call,grouping\|\|Left to right\|\|Highest -1\| \|\|\|Array access\|\|Left to right\|\|1\| \|.\|\|Object property access\|\|Left to right\|\|1\| \|++\|\|Increment\|\|Right to left\|\|2\| \|--\|\|Decrement\|\|Right to left\|\|2\| \|-\|\|Negation\|\|Right to left\|\|2\| \|!\|\|Logical NOT\|\|Right to left\|\|2\| \|-\|\|Bitwise NOT\|\|Right to left\|\|2\| \|delete\|\|Removes array value or object property\|\|Right to left\|\|2\| \|new\|\|Creates an object\|\|Right to left\|\|2\| \|typeof\|\|Returns data type\|\|Right to left\|\|2\| \|void\|\|Specifies no value to return\|\|Right to left\|\|2\| \|/\|\|Division\|\|Left to right\|\|3\| \|\|\|Multiplication\|\|Left to right\|\|3\| \|%\|\|Modulus\|\|Left to right\|\|3\| \|+\|\|Plus\|\|Left to right\|\|4\| \|+\|\|String Concatenation\|\|Left to right\|\|4\| \|-\|\|Subtraction\|\|Left to right\|\|4\| \|»\|\|Bitwise right-shift\|\|Left to right\|\|5\| \|«\|\|Bitwise left-shift\|\|Left to right\|\|5\| \|>.>=\|\|Greater than, greater than or equal to\|\|Left to right\|\|6\| \|<,<=\|\|Less than, less than or equal to\|\|Left to right\|\|6\| \|==\|\|Equality\|\|Left to right\|\|7\| \|!=\|\|Inequality\|\|Left to right\|\|7\| \|===\|\|Identity operator - equal to (and same data type)\|\|Left to right\|\|7\| \|!==\|\|Non-identity operator - not equal to (or don't have the same data type)\|\|Left to right\|\|7\| \|&\|\|Bitwise AND\|\|Left to right\|\|8\| \|^\|\|Bitwise XOR\|\|Left to right\|\|9\| \|\|\|\|Bitwise OR\|\|Left to right\|\|10\| \|&&\|\|Logical AND\|\|Left to right\|\|11\| \|Il\|\|Logical OR\|\|Left to right\|\|12\| \|?:\|\|Conditional branch\|\|Left to right\|\|13\| \|=\|\|Assignment\|\|Right to left\|\|14\| \|*=,/=,%=,+=,,==,<<=,>>=, >>>=, &=,^=,1=\|\|Assignment according to the preceding operator\|\|Right to left\|\|14\| Now you can use the above table to validate all the examples demonstrated previously. We have Bitwise XOR, the modulus operator, and the exponentiation operator as well. Here's how we'll evaluate it: We evaluate the expression from left to right and carry out the XOR operation first. Then, we evaluate the exponentiation operation because of right to left associativity. Then we evaluate the modulus operation because the modulus operator has higher precedence than subtraction which gives us the result 20. Grouping and Short Circuiting If you look at the table again, you'll notice that the grouping operator () has the highest precedence. This means that the given expression will evaluate to 16 and not 13: Note that while demonstrating which operator will be evaluated first, I have manually applied the grouping operator in all the previous examples. That was simply an attempt to distinguish the operation being performed at that instance from the rest of the operations present in the expression. If the grouping operator has the highest precedence, what do you think will happen in the following statement: Here, the above expression will evaluate false. Well, you might think that since we have a grouping operator, the addition operation is evaluated first, and then the conditional or logical AND operator is evaluated. However, that's not the case here. This is due to a concept called short-circuiting. In the above statement, we know false is not true. So the statement or expression is short-circuited at that point, and the grouping operator is not evaluated at all. - The grouping operator is assigned the highest level of precedence. - Short-circuiting allows conditional operators to take precedence over the grouping operator.
About This Chapter Studying for Physics 112 - Summary The flashcard sets contained within this collection are here to help you reinforce your understanding of Physics II course material. As you flip through the cards in each set, you'll hone your recognition of formulas, definitions, concepts, laws and theories, famous figures and more. Use these cards any time you want to improve your knowledge of physics, whether you are a science student or just interested in the subject; it's a great way to study for a test or exam, too! The flashcard sets in this chapter cover such topics as: - Energy and heat transfer - Thermodynamics, the ideal gas law, and kinetic theory - Light, optics and electromagnetic waves - Electrostatics and magnetism - Circuits, capacitors and electronics - Relativity and quantum theory How They Work It's so easy to use our Physics 112 flashcards! Master those tricky formulas, definitions, laws and concepts by following these simple instructions: - Open a flashcard set for the physics topic you're studying. - Identify the keyword, formula, image, definition or other physics concept on the flashcard. - Click the button to flip the card and view the answer. - Select whether you got it right or missed it to keep track of your progress. - Click the Next arrow to proceed through the flashcard set. - Review the cards you missed or shuffle all cards back in and go through them again. Energy Transfers in Physics Flashcards In this set of flashcards, you will review how energy is transferred in physical systems. These flashcards will give a basic overview of energy transfers and then will focus in on reactions involving heat and electrical energy. Heat Transfer & Interaction Flashcards How does a hot cup of coffee left on the table cool off over time? These flashcards will help you review how heat moves from one object to another and how to complete calculations involving heat energy transfer. Basic Nuclear Physics Flashcards Use these flashcards to review basic nuclear physics, including the particles that make up an atom's nucleus (and the particles that make up those particles) and the physical phenomena that they are part of. 4. Physics 112: Physics II Formulas & Constants The Physics 112: Physics II course includes formulas related to energy, electric current, electromagnet waves and magnetic fields, most of which you'll find on this sheet. Earning College Credit Did you know… We have over 200 college courses that prepare you to earn credit by exam that is accepted by over 1,500 colleges and universities. You can test out of the first two years of college and save thousands off your degree. Anyone can earn credit-by-exam regardless of age or education level. To learn more, visit our Earning Credit Page Transferring credit to the school of your choice Not sure what college you want to attend yet? Study.com has thousands of articles about every imaginable degree, area of study and career path that can help you find the school that's right for you. Other chapters within the Physics 112: Physics II course - Energy Transfers in Physics - Heat Transfer & Interaction - The Basics of Thermodynamics - Ideal Gas Law & Kinetic Theory - Light & Electromagnetic Waves - Mirrors & Lenses in Geometric Optics - Basics of Electrostatics - Magnetism Basics - Series, Parallel & Combined Circuits - Capacitors, Inductors & Alternating Current - Modern Quantum Theory - Basic Nuclear Physics - General & Special Relativity
Huge product rangeOver 140,000 books & equipment products Rapid shippingUK & Worldwide Pay in £, € or U.S.$By card, cheque, transfer, draft Exceptional customer serviceGet specialist help and advice This user-friendly interactive book makes concepts simple and gives you the confidence and knowledge to learn and teach mathematics to primary children. You can conveniently access easily digestible content to build and test your knowledge. The SAGE Primary Mathematics Student Panel includes: - Lesson plans and worksheets: Save time with ideas and resources for planning your lessons. - Problem solved! videos: Visualise how to solve a question by watching a quick animated demonstration. - Learning and teaching points: Feel confident in the classroom with these key points to remember when planning and teaching. - Questions and quizzes: Achieve a firm grasp of concepts and a clear understanding with self-assessment questions and quizzes. - Free interactive ebook: Study anywhere with your portable and convenient eBook! SECTION A: MATHEMATICAL UNDERSTANDING Chapter 1: Primary Teachers’ Insecurity about Mathematics Chapter 2: Mathematics in the Primary Curriculum Chapter 3: Learning how to Learn Mathematics SECTION B: MATHEMATICAL REASONING AND PROBLEM SOLVING Chapter 4: Key Processes in Mathematical Reasoning Chapter 5: Modelling and Problem Solving SECTION C: NUMBERS AND CALCULATIONS Chapter 6: Numbers and Place Value Chapter 7: Addition and Subtraction Structures Chapter 8: Mental Strategies for Addition and Subtraction Chapter 9: Written Methods for Addition and Subtraction Chapter 10: Multiplication and Division Structures Chapter 11: Mental Strategies for Multiplication and Division Chapter 12: Written Methods for Multiplication and Division SECTION D: FURTHER NUMBER CONCEPTS AND SKILLS Chapter 13: Natural Numbers: Some Key Concepts Chapter 14: Integers: Positive and Negative Chapter 15: Fractions and Ratios Chapter 16: Decimal Numbers and Rounding Chapter 17: Calculations with Decimals Chapter 18: Proportionality and Percentages SECTION E: ALGEBRA Chapter 19.: Algebraic Reasoning Chapter 20: Coordinates and Linear Relationships SECTION F: MEASUREMENT Chapter 21: Concepts and Principles of Measurement Chapter 22: Perimeter, Area and Volume SECTION G: GEOMETRY Chapter 23: Angle Chapter 24: Transformations and Symmetry Chapter 25: Classifying Shapes SECTION H: STATISTICS AND PROBABILITY Chapter 26: Handling Data Chapter 27: Comparing Sets of Data Chapter 28: Probability Derek Haylock is an education writer, with an extensive list of publications in the field of mathematics education. He worked for over 30 years in teacher education, both initial and in-service, and was Co-Director of Primary Initial Teacher Training and responsible for the mathematics components of the primary programmes at the University of East Anglia (UEA) in Norwich. He has considerable practical experience of teaching and researching in primary classrooms. His work in mathematics education has taken him to Germany, Belgium, Lesotho, Kenya, Brunei, India and Sweden. As well as his publications in the field of education, he has written seven books of Christian drama for young people and a Christmas musical (published by Church House/National Society). For 15 years after his time at UEA, he was in great demand as a consultant and professional speaker. His work as a writer continues. Ralph Manning is an independent consultant in primary education. He has worked as a primary teacher and as a lecturer in primary teacher education for 18 years, following an earlier career in the IT industry. Ralph led the mathematics and physical education programmes for the Primary PGCE at the University of East Anglia until 2013. His other specific interests are in developing children's higher thinking skills, assessment for learning and planning for teaching effectively. He was a former founding member of the General Teaching Council for England and is engaged in supporting developing education projects in India.
Definition of numeracy Numeracy is the capacity to take mathematics and apply knowledge, skills and strategies to deal with everyday life in a variety of situations. It will also give students the ability to cope confidently with the mathematical demands of further education, employment and adult life. Kingsthorpe College is committed to raising the standards of Numeracy of all its students, so that they develop the ability to use Numeracy skills in all areas of the curriculum and the skills necessary to cope confidently with the demands of further education, employment and adult life. A numerate student is one who - Is confident and competent in performing calculations involving number - Can use a range of techniques to carry out computations mentally and on paper - Knows and understands the properties of number - Can explain methods and justify reasoning and conclusions, using correct mathematical terms - Can use calculators and other ICT resources appropriately and effectively to solve mathematical problems and select from the display the number of figures appropriate to the context of the calculation - Can recognise and use mathematical skills and techniques in a variety of contexts Parents can support their children by - Referring to the student’s exercise book in order to see the topics currently being worked upon - Ensuring that emphasis is placed on mental and pencil and paper methods for calculations and that calculators are only used where appropriate when working at home (method and working must be shown irrespective of calculator use) - Valuing their child’s own methods of doing calculations - Purchase of a revision guide appropriate to their children’s level / Grade (available in College or from bookshops) - Encouraging their children to make use of the learning tools in the following websites: www.mymaths.co.uk and www.samlearning.com Numeracy across the curriculumScience Almost every scientific investigation or experiment is likely to require one or more of the mathematics skills of classifying, counting, measuring, calculating, estimating and recording in tables and graphs.Art, Design and Technology Measurements are often needed in art, design and technology. Many patterns and constructions in our own and other cultures are based on spatial ideas and properties of shape, including symmetry. Designs may need enlarging or reducing, introducing ideas of multiplication, scale and ratio. The preparation of food involves measurement, working out times and calculating cost, frequently extending into calculations involving ratio and proportion.ICT In ICT lessons, students will collect and classify data, enter values into data handling software, produce graphs and tables and interpret and explain their results. Spreadsheets require algebraic and graphical skills involving constructing formulae and generating sequences, functions and graphs.Humanities Students will make statistical enquiries involving primary and secondary data and the interpretation of graphs, charts and tables. Students will apply mathematical skills to financial problems and other real-life contexts such as the study of maps involving co-ordinates, angle, direction, position, scale and ratio.PE and Music Athletic activities use measurement of height, distance and time, data logging devices to quantify, explore and improve performance. Ideas of counting, time, symmetry, movement, position and direction are used extensively in music, drama, gymnastics, athletics and competitive games.Year 7 catch up premium and numeracy Kingsthorpe College received an additional premium for each Year 7 pupil who did not achieve at least a level 4 in reading and/or maths at Key Stage 2. These additional funds have been used to deliver additional tuition and intensive support for students in small groups, giving students valuable support to bring them up to speed to ensure that they are more likely to succeed. The school will provide - Individual tuition in addition to classroom teaching. - Intensive small group tuition supported by pertinent materials and resources which include after school sessions. - One to one session particularly for assessment. - Small group work for numeracy. - Group work to identify skills gaps for students followed by tailored sessions to facilitate personalised learning.
The Parliament of the United Kingdom of Great Britain and Northern Ireland is the supreme legislative institution in the United Kingdom and British overseas territories (it alone has parliamentary sovereignty). At its head is the Sovereign; it also includes an Upper House, called the House of Lords, and a Lower House, called the House of Commons. The House of Lords includes two different types of members—the Lords Spiritual (the senior clergy of the Church of England) and the Lords Temporal (members of the Peerage); it is an almost wholly appointed body. The House of Commons, on the other hand, is a democratically elected chamber. The House of Lords and the House of Commons meet in separate chambers in the Palace of Westminster (commonly known as the "Houses of Parliament"), in British capital, London (more exactly, in the borough known as the City of Westminster). Parliament evolved from the ancient council which advised the Sovereign. In theory, power is vested not in Parliament, but in the "Queen-in-Parliament" (or "King-in-Parliament"). The Queen-in-Parliament is often said to be a completely sovereign authority, though such a position is debatable. In modern times, real power is vested in the democratically elected House of Commons; the Sovereign acts only as a figurehead, and the powers of the House of Lords are greatly limited. The British Parliament is often called the "Mother of Parliaments," as the legislative bodies of many nations—most notably, those of the members of the Commonwealth—are modeled on it. However, it is a misquotation of John Bright, who had actually remarked on 18 January 1865 that "England is the Mother of Parliaments", in the context of supporting demands for expanded voting rights in a country which had pioneered Parliamentary government. One may trace the origin of Parliament to the times of the Anglo-Saxons. Anglo-Saxon Kings were advised by a council known as the Witenagemot, whose foremost members were the King's sons and brothers. The Earldormen, or executive heads of the shires, also had seats in the Witenagemot, as did the senior clergymen of the state. The King still possessed ultimate authority, but laws were made only after seeking the advice (and, in later times, the consent) of the Witenagemot. The entire Anglo-Saxon body politic was reformed when William of Normandy conquered England in 1066. William brought to England the feudal system he was accustomed to in his native France. Thus, he granted land to his most important military supporters, who in turn granted land to their supporters, thus creating a feudal hierarchy. Those who held lands directly from the King were known as tenants-in-chief, and the territories they held were called manors. William I was an absolute ruler, but, as a matter of course, he sought the advice of a council of tenants-in-chief and ecclesiastics, before making laws. The tenants-in-chief often struggled with their spiritual counterparts and with the King for power. In 1215, they secured from John the Magna Carta, which established that the King may not levy or collect any taxes (except the feudal taxes to which they were hitherto accustomed), save with the consent of his council. It was also established that the most important tenants-in-chief (the earls and the barons), as well as the ecclesiastics (archbishops, bishops and abbots) be summoned to the council by personal writs from the Sovereign, and that all others be summoned to the council by general writs from the sheriffs of their counties. John later repealed the Magna Carta, but Henry III reinstated it. The royal council slowly developed into a Parliament. In 1265, Simon de Montfort, 6th Earl of Leicester, who was in rebellion against King Edward I, summoned a parliament of his supporters without any or prior royal authorisation. The archbishops, bishops, abbots, earls and barons were summoned, as were two knights from each shire and two burgesses from each borough. Knights had been summoned to previous councils, but the representation of the boroughs was unprecedented. De Montfort's scheme was formally adopted by Edward I in the so-called "Model Parliament" of 1295. At first, each estate debated independently; by the reign of Edward III, however, Parliament had been separated into two Houses: one, including the nobility and higher clergy, the other, including the knights and burgesses. The authority of Parliament grew under Edward III; it was established that no law could be made, nor any tax levied, without the consent of both Houses as well as of the Sovereign. The growing influence of Parliament was restrained by numerous civil wars. By the end of the Wars of the Roses, royal supremacy had been restored. The Crown was at the height of its power during the reign of Henry VIII. The number of the Lords Spiritual diminished under Henry, who commanded the Dissolution of the Monasteries, thereby depriving the abbots and priors of their seats in the Upper House. For the first time, the Lords Temporal were more numerous than the Lords Spiritual. Parliaments continued to behave submissively under the Tudor monarchs who followed Henry, but began to display an unusual sense of independence under Elizabeth I. As England evolved into a world power, members of both Houses actively discussed succession to the Crown (the Queen never married) and condemned various royal policies. Their new-found boldness proved intolerable to Elizabeth's Scottish successor, James I (who was simultaneously King in Scotland as James VI). The great struggle between the Crown and Parliament occurred under James I's successor, Charles I. Alarmed by the arbitrary exercise of royal power, the House of Commons submitted to Charles the Petition of Right, demanding the restoration of their liberties, in 1628. Though he accepted the petition, he later dissolved Parliament and ruled without them for eleven years. It was only after the financial disaster of the Scottish Bishops' Wars (1639–1640) that he was forced to recall Parliament in order that they may authorise new taxes. The new Parliament was quite rebellious; their struggle for power with the Crown culminated in the English Civil War. In 1649, Charles was executed and replaced by the military dictator Oliver Cromwell. The House of Lords was abolished, and the House of Commons remained subordinate to Cromwell. After Cromwell's death, however, the monarchy was restored in 1660. The House of Lords also returned. Following the Restoration, monarchs undertook to regularly summon Parliament. Nevertheless, there was no explicit guarantee of Parliamentary liberties until James II, an unpopular Catholic ruler, was forced to flee the country in 1688. Parliament "deemed" that he had abdicated, but it offered the Crown to his Protestant daughter Mary, instead of his Catholic son. Mary II ruled jointly with her husband, William III. They agreed to the Act Declaring the Rights and Liberties of the Subject and Settling the Succession of the Crown (the Bill of Rights), which acknowledged several powers Parliaments had claimed under previous Stuart monarchs. No taxes were to be levied, nor any standing armies to be kept during peacetime, without the consent of Parliament. Freedom of speech in Parliamentary debates was also secured. The influence of Parliament was undoubtedly augmented by the Bill of Rights, but the supremacy of the Crown stil remained clear. It was only after the Hanoverian George I ascended the Throne in 1714 that power began to shift from the Sovereign. George was a German ruler, spoke poor English and preferred to concentrate on his dominions in Europe. He thus entrusted power to a group of his ministers, the foremost of which was Sir Robert Walpole. George III sought to restore royal supremacy, but by the end of his reign, the position of the ministers—who would in turn have to rely on Parliament for support—was cemented. The principle of ministerial responsibility to the Lower House did not develop until the nineteenth century. The House of Lords was superior to the House of Commons both in theory and in practice. Members of the House of Commons were elected in an antiquated electoral system, under which constituencies of vastly different sizes existed. Thus, the borough of Old Sarum, with seven voters, could elect two members, as could the borough of Dunwich, which had completely disappeared into the sea due to land erosion. In many cases, members of the Upper House controlled tiny constituencies, known as pocket boroughs or rotten boroughs, and could ensure the election of their relatives or supporters. Many seats in the House of Commons were "owned" by the Lords. After the reforms of the nineteenth century (beginning in 1832), the electoral system in the Lower House was much more regularised. No longer dependent on the Upper House for their seats, members of the House of Commons began to grow more assertive. The supremacy of the House of Commons was clearly established during the early twentieth century. In 1909, the Commons passed the so-called "People's Budget," which made numerous changes to the taxation stystem in a manner detrimental to wealthy landowners. The House of Lords, which mostly consisted of powerful landowning aristocrats, rejected the Budget. On the basis of the Budget's popularity and the Lords' consequent unpopularity, the Liberal Party won a general election in 1910. Using the result as a mandate, the Liberal Prime Minister, Herbert Henry Asquith, introduced the Parliament Bill, which sought to restrict the powers of the House of Lords. (He did not reintroduce the land tax provision of the People's Budget.) When the Lords refused to pass the bill, Asquith approached the King and requested the creation of several hundred Liberal peers so as to erase the Conservative majority in the House of Lords. In the face of such a threat, the House of Lords reluctantly passed the bill. The Parliament Act 1911, as it became known, allowed the Lords to delay a bill for a maximum of three sessions (reduced to two sessions in 1949), after which it could become law over their objections. Further reforms to the House of Lords have been made during the twentieth century. In 1958, the Life Peerages Act authorised the regular creation of life peerage dignities. By the 1960s, the regular creation of hereditary peerage dignities had ceased; thereafter, almost all new peers were life peers only. More recently, the House of Lords Act 1999 has removed the automatic right of hereditary peers to sit in the Upper House (although it made an exception for ninety-two of them on a temporary basis). The House of Lords is now a chamber that is, in practice, subordinate to the House of Commons. At the head of Parliament is the British Sovereign. The Sovereign's role, however, is merely ceremonial; in practice, he or she always acts on the advice of the Prime Minister and other ministers, who are in turn accountable to the two Houses of Parliament. The Upper House, the House of Lords, is mostly made up of appointed members ("Lords of Parliament"). The whole House is formally styled The Right Honourable The Lords Spiritual and Temporal in Parliament Assembled, the Lords Spiritual being clergymen of the Church of England and the Lords Temporal being Peers of the Realm. The Lords Spiritual and Lords Temporal are considered separate "estates," but they sit, debate and vote together. The Lords Spiritual formerly included all of the senior clergymen of the Church of England—archbishops, bishops, abbots and priors. Upon the Dissolution of the Monasteries under Henry VIII, however, the abbots and priors lost their positions in Parliament. All diocesan bishops continued to sit in Parliament, but the Bishopric of Manchester Act 1847 and later acts provide that only the twenty-six most senior are Lords Spiritual. These twenty-six always include the incumbents of the "five great sees," namely, the Archbishop of Canterbury, the Archbishop of York, the Bishop of London, the Bishop of Durham and the Bishop of Winchester. The remaining twenty-one Lords Spiritual are the most senior diocesan bishops, ranked in order of consecration. The Lords Temporal are all members of the Peerage. Formerly, they included hereditary peers, of the ranks of Duke, Marquess, Earl, Viscount and Baron. The right of some hereditary peers to sit in Parliament was not automatic: after Scotland and England united into Great Britain in 1707, it was provided that all peers whose dignities had been created by English Kings could sit in Parliament, but those whose dignities had been created by Scottish Kings were to elect a limited number of "representative peers." A similar arrangement was made in respect of Ireland when that nation merged with Great Britain in 1801. But when Southern Ireland left the United Kingdom in 1922, the election of Irish representative peers ceased. By the Peerage Act 1963, the election of Scottish representative peers also ended, but all Scottish peers were granted the right to sit in Parliament. Under the House of Lords Act 1999, only life peerage dignities (that is to say, peerage dignities which cannot be inherited) automatically entitle their holders to seats in the House of Lords. Of the hereditary peers, only ninety-two—the individuals exercising the offices of Earl Marshal and Lord Great Chamberlain, in addition to ninety hereditary peers elected by other peers—retain their seats in the House. The Commons, the last of the "estates" of the Kingdom, are represented in the House of Commons, which is formally styled The Honourable The Commons in Parliament Assembled. The House consists of 659 members ("Members of Parliament" or "MPs"), each of whom is chosen by a single constituency according to the First-Past-the-Post electoral system. Universal adult suffrage prevails; citizens of the United Kingdom, as well as citizens of the Republic of Ireland and of Commonwealth nations resident in the United Kingdom, are qualified to vote. The term of members of the House of Commons depends on the term of Parliament; a general election, during which all the seats are contested, occurs after each dissolution (see below). The three branches of Parliament are supposed to be kept separate from each other; no individual may form a part of more than one component of Parliament. Lords of Parliament are legally barred from voting in elections for members of the House of Commons; furthermore, the Sovereign by convention does not vote, although there is no statutory impediment. Each of the two Houses of Parliament is presided over by a Speaker. In the House of Lords, the Lord Chancellor, a member of the Cabinet, is the ex officio Speaker. Where there is a vacancy in the office, a Speaker may be appointed by the Crown. Deputy Speakers, who take the place of an absent Lord Chancellor, are also chosen by the Crown. The House of Commons has the right to elect its own Speaker. Theoretically, the approval of the Sovereign is required before the election becomes valid, but it is, by modern conventions, always granted. The Speaker's place may be taken by three deputies, known as the Chairman, First Deputy Chairman and Second Deputy Chairman of Ways and Means. (They take their name from the Committee of Ways and Means, of which they were once presiding officers, but which no longer exists.) In general, the Lord Chancellor's influence as Speaker is very limited, whilst the powers belonging to the Speaker of the House of Commons are vast. Decisions on points of order and on the disciplining of unruly members are made by the whole body in the Upper House, but by the Speaker alone in the Lower House. Speeches in the House of Lords are addressed to the House as a whole (using the words "My Lords"), but those in the House of Commons are addressed to the Speaker alone (using the words "Mr Speaker" or "Madam Speaker"). Both Houses may decide questions with voice voting; members shout out "Aye" and "No" (in the House of Commons), or "Content" and "Not-Content" (in the House of Lords), and the presiding officer declares the result. The pronouncement of the Lord Chancellor or Speaker may be challenged, and a recorded vote (known as a division) demanded. (The Speaker of the House of Commons may choose to overrule a frivolous request for a division, but the Lord Chancellor does not possess an equivalent power.) In each House, a division requires members to file into one of the two lobbies alongside the Chamber; their names are recorded by clerks, and their votes are counted as they exit the lobbies to re-enter the Chamber. The Speaker of the House of Commons, who is expected to remain non-partisan, does not cast a vote except in the case of a tie; the Lord Chancellor, however, votes along with the other Lords. (For further details on procedure, see the separate articles on the House of Lords and the House of Commons.) Following a general election, a new Parliamentary session begins. Parliament is formally summoned forty days in advance by the Sovereign, who is considered the source of parliamentary authority. On the day indicated by the Sovereign's proclamation, the two Houses assemble in their respective chambers. The Commons are then summoned to the House of Lords, where Lords Commissioners (representatives of the Sovereign) instruct them to elect a Speaker. The Commons perform the election; on the next day, they return to the House of Lords, where the Lords Commissioners confirm the election and grant the new Speaker the royal approval in the Sovereign's name. The business of Parliament for the next few days of its session involves the taking of the oaths of allegiance. Once a majority of the members have taken the oath in each House, the State Opening of Parliament may occur. The Lords take their seats in the House of Lords Chamber, the Commons appear at the Bar (immediately outside the Chamber), and the Sovereign takes his or her seat on a throne. The Sovereign then reads the Speech from the Throne—the content of which is determined by the Ministers of the Crown—outlining the Government's legislative agenda for the upcoming year. Thereafter, each House proceeds to the transaction of legislative business. By custom, before considering the Government's legislative agenda, a bill is introduced pro forma in each House—the Select Vestries Bill in the House of Lords and the Outlawries Bill in the House of Commons. These bills do not actually become laws; they are merely ceremonial indications of the power of each House to debate independently of the Crown. After the pro forma bill is introduced, each House debates the content of the Speech from the Throne for several days. Once each House formally sends its reply to the Speech, the proper legislative business of the House may commence. At once, each House becomes fully active in appointing committees, electing officers, passing resolutions and considering legislation. A session of Parliament is brought to an end by a prorogation. There is a ceremony similar to the State Opening, but it is much less well-known. Normally, the Sovereign does not personally attend the prorogation ceremony in the House of Lords; rather, he or she is represented by Lords Commissioners. The next session of Parliament begins under the procedures described above, but it is not necessary to conduct another election of a Speaker or take the oaths of allegiance afresh at the beginning of such subsequent sessions. Instead, the State Opening of Parliament is proceeded to directly. Each Parliament, after a number of sessions, comes to an end, either by the command of the Sovereign or by effluxion of time, the former being more common in modern times. The dissolution of Parliament is effected by the Sovereign, but always on the advice of the Prime Minister. The Prime Minister may seek a dissolution because the time is politically advantageous to his or her party. Furthermore, if the Prime Minister loses the support of the House of Commons, he must either resign or seek a dissolution of Parliament to renew his or her mandate. Originally, there was no fixed limit on the length of a Parliament, but the Triennial Act 1694 set the maximum duration at three years. As the frequent elections were deemed inconvenient, the Septennial Act 1716 extended the maximum duration to seven years, but the Parliament Act 1911 reduced it to five years. During the Second World War, the term was temporarily extended to ten years by Acts of Parliament. Since the end of the war in 1945, however, the maximum term has remained five years. Modern Parliaments, however, rarely continue for the maximum duration; normally, they are dissolved earlier. For instance, the Fifty-Second Parliament assembled in 1997, but was dissolved after only four years. Formerly, the demise of the Sovereign automatically brought a Parliament to an end, for the Crown was seen as the caput, principium, et finis (beginning, basis and end) of the body. It was, however, deemed inconvenient to have no Parliament at a time when succession to the Crown could be disputed. Thus, a statute passed during the reign of William III and Mary II provided that a Parliament was to continue for six months after the death of a Sovereign, unless dissolved earlier. The Representation of the People Act 1867 brought this arrangement to an end; now, a demise in the Crown does not affect the duration of a Parliament. After each Parliament concludes, a general election is held, and new members of the House of Commons elected. The membership of the House of Lords, however, does not change due to a dissolution. Each Parliament which assembles following a general election is deemed to be distinct from the one which just concluded. Thus, each Parliament is separately numbered, the present Parliament being the Fifty-Third Parliament of the United Kingdom (that is to say, the fifty-third Parliament summoned since the formation of the United Kingdom of Great Britain and Ireland in 1801). Previous Parliaments were "of Great Britain" or "of England." Parliament meets in the Palace of Westminster. Laws, in draft form known as bills, may be introduced by any member of either House. Usually, however, a bill is introduced by a Minister of the Crown. A bill introduced by a Minister is known as a "Government Bill"; one introduced by another member is called a "Private Member's Bill." A different way of categorising bills involves the subject. Most bills, involving the general public, are called "Public Bills." A bill that seeks to grant special rights to an individual or small group of individuals is called a "Private Bill." A Private Bill which has broader public implications is called a "Hybrid Bill." Each Bill goes through several stages in each House. The first stage, called the first reading, is a mere formality. At the next stage, the second reading, the general principles of the bill are debated. At the second reading, the House may vote to reject the bill (by refusing to pass the motion "That the Bill be now read a second time"), but defeats of Government Bills are extremely rare. Following the second reading, the bill is sent to a committee. In the House of Lords, the Committee of the Whole House or the Grand Committee is used. Each consists of all members of the House; the latter operates under special procedures, and is used only for uncontroversial bills. In the House of Commons, the bill is usually committed to a Standing Committee, consisting of between sixteen and fifty members, but the Committee of the Whole House is used for important legislation. Several other types of committees, including Select Committes, may be used, but are in practice only rarely employed. A committee considers the bill clause-by-clause, and reports its proposed amendments to the entire House, where further detailed consideration occurs. Once the House considers the bill, the third reading follows. In the House of Commons, no further amendments may be made, and the passage of the motion "That the Bill be now read a third time" amounts to passage of the whole bill. In the House of Lords, however, further amendments to the bill may be moved. After the passage of the third reading motion mentioned above, the House of Lords must vote on another motion "That the Bill do now pass." Following its passage in one House, the bill is sent to the other House. If passed in identical form by both Houses, it may be presented for the Sovereign's Assent. If, however, one House passes amendments that the other will not agree to, and the two Houses cannot resolve their disagreements, the bill fails. Since the passage of the Parliament Act 1911, however, the power of the House of Lords to interfere with bills passed by the House of Commons has been restricted. Further restrictions were placed by the Parliament Act 1949. Thus, if the House of Commons passes a public bill in two successive sessions, and the House of Lords rejects them both times, then the Commons may direct that the bill be presented to the Sovereign for his or her Assent, disregarding the rejection of the Bill in the House of Lords. In each case, the bill must be passed by the House of Commons at least one calendar month before the end of the session. The provision does not apply to bills originated in the House of Lords, to bills seeking to extend the duration of a Parliament beyond five years or to Private Bills. A special procedure applies in relation to bills classified by the Speaker of the House of Commons as "Money Bills." A Money Bill solely concerns national taxation or public funds; the Speaker's certificate is deemed conclusive under all circumstances. If the House of Lords fails to pass a Money Bill within one month of its passage in the House of Commons, the Lower House may direct that the Bill be submitted for the Sovereign's Assent immediately. Even before the passage of the Parliament Acts, the Commons possessed pre-eminence in cases of financial matters. By ancient custom, the House of Lords may neither introduce a bill relating to taxation or Supply, nor amend a bill so as to insert a provision relating to taxation or Supply, nor amend a Supply Bill in any way. The House of Commons, however, is free to waive this privilege, and sometimes does so to allow the House of Lords to pass amendments with financial implications. The House of Lords, however, remains free to reject bills relating to Supply and taxation, but may be easily overruled if the bills are Money Bills. (A bill relating to revenue and Supply may not be a Money Bill if, for example, it includes subjects other than national taxation and public funds). The last stage of a bill involves the granting of the Royal Assent. Theoretically, the Sovereign may grant the Royal Assent (that is, make the bill a law) or withhold the Royal Assent (that is, veto the bill). Under modern notions of a constitutional monarchy, however, the Sovereign always grants the Royal Assent. The last refusal to grant the Assent came in 1708, when Anne withheld her Assent from a bill "for the settling of Militia in Scotland." Every bill, thus, obtains the assent of all three components of Parliament before it becomes law (except as provided by the Parliament Acts where the House of Lords is overriden). All laws are in theory "enacted" by the Sovereign, with the consent of the Lords and Commons. The words "BE IT ENACTED by the Queen's [King's] most Excellent Majesty, by and with the advice and consent of the Lords Spiritual and Temporal, and Commons, in this present Parliament assembled, and by the authority of the same, as follows:-" form a part of each Act of Parliament (where the House of Lords' authority has been overriden through the usage of the Parliament Acts, the words "BE IT ENACTED by The Queen's [King's] most Excellent Majesty, by and with the advice and consent of the Commons in this present Parliament assembled, in accordance with the provisions of the Parliament Acts 1911 and 1949, and by the authority of the same, as follows:-" are used instead). These words at the beginning of every Act is known as the enacting formula. In addition to its legislative functions, Parliament also performs several judicial functions. The Queen-in-Parliament constitutes the highest court in the realm for most purposes, but the Privy Council has jurisdiction in some cases (for instance, appeals from ecclesiastical courts). The jurisdiction of Parliament arises from the ancient custom of petitioning the Houses to redress grievances and to do justice. The House of Commons ceased considering petitions to reverse the judgements of lower courts in 1399, effectively leaving the House of Lords as the realm's court of last resort. In modern times, the judicial functions of the House of Lords are performed not by the whole House, but by a group of "Lords of Appeal in Ordinary" (judges granted life peerage dignities under the Appellate Jurisdiction Act 1876 by the Sovereign) and by "Lords of Appeal" (other peers with experience in the judiciary). The Lords of Appeal in Ordinary and Lords of Appeal (or "Law Lords") are Lords of Parliament, but normally do not vote or speak on political matters. Certain other judicial functions have historically been performed by the House of Lords. Until 1948, it was the body in which peers of the Realm had to be tried for felonies or high treason; now, peers are tried by normal juries. Furthermore, when the House of Commons impeaches an individual, the trial takes place in the House of Lords. Impeachments, however, are now obsolete; the last impeachment occurred in 1806. Relationship with the Government The British Government is answerable to the House of Commons. However, neither the Prime Minister nor members of the Government are elected by the House of Commons. Instead, the Queen requests the person most likely to command the support of a majority in the House, normally the leader of the largest party in the House of Commons, to form a government. So that they may be accountable to the Lower House, the Prime Minister and most members of the Cabinet are members of the House of Commons instead of the House of Lords. The last Prime Minister to be a Lord of Parliament was Alec Douglas-Home, 14th Earl of Home, who became Prime Minister in 1963. Nevertheless, to adhere to the convention under which he was responsible to the Lower House, Lord Home disclaimed his peerage dignity and procured election to the House of Commons. Parliament controls the executive by passing or rejecting its Bills and by forcing Ministers of the Crown to answer for their actions, either at "Question Time" or during meetings of the parliamentary committees. In both cases, the Ministers are asked questions by members of their Houses, and are obliged to answer. Although the House of Lords may scrutinise the executive through Question Time and through its committees, it cannot bring about the end of a Government. A ministry must, however, always retain the confidence and support of the House of Commons. The Lower House may indicate its lack of support by rejecting a Motion of Confidence or by passing a Motion of No Confidence. Confidence Motions are generally originated by the Government in order to reinforce its support in the House, whilst No Confidence Motions are introduced by the Opposition. The motions sometimes take the form "That this House has [no] confidence in Her Majesty's Government" but several other varieties, many referring to specific policies supported or opposed by Parliament, are often used. For instance, a Confidence Motion of 1992 used the form, "That this House expresses the support for the economic policy of Her Majesty's Government." Such a motion may theoretically be introduced in the House of Lords, but, as the Government need not enjoy the confidence of that House, would not be of the same effect as a similar motion in the House of Commons; the only modern instance of such an occurrence involves the No Confidence Motion that was introduced in 1993, and subsequently defeated. Many votes are considered votes of confidence, although not specifically involving the language mentioned above. Important bills that form part of the Government's agenda (as stated in the Speech from the Throne) are generally considered matters of confidence. The defeat of such a bill by the House of Commons indicates that a Government no longer has the confidence of that House. Furthermore, the same effect is achieved if the House of Commons "withdraws Supply," that is, rejects the Budget. Where a Government has lost the confidence of the House of Commons, the Prime Minister is obliged to either resign, or seek the dissolution of Parliament and a new general election. Where a Prime Minister has ceased to retain a majority in that vote and requests a dissolution, the Sovereign can in theory reject his request, forcing his resignation and allowing the Leader of the Opposition to be asked to form a new government. This power however is supposed to be used extremely rarely. The conditions that should be met to allow such a refusal are known as the Lascelles Principles. Note, however, that these conditions and principles are merely informal conventions; it is possible, though highly improbable, for the Sovereign to refuse a dissolution for no reason at all. In practice, the House of Commons' scrutiny of the Government is very weak. Since the First-Past-the-Post electoral system is employed in elections, the governing party tends to enjoy a large majority in the Commons; there is often limited need to compromise with other parties. Modern British political parties are so tightly organised that they leave relatively little room for free action by their MPs. In many cases, MPs may be expelled from their parties for voting against the instructions of party leaders. During the twentieth century, the Government has lost confidence issues only thrice—twice in 1924, and once in 1979. Several different views have been taken of Parliament's sovereignty. According to the jurist Sir William Blackstone, "It has sovereign and uncontrollable authority in making, confirming, enlarging, restraining, abrogating, repealing, reviving, and expounding of laws, concerning matters of all possible denominations, ecclesiastical, or temporal, civil, military, maritime, or criminal … It can, in short, do every thing that is not naturally impossible." A different view, however, has been taken by the Scottish judge Lord Cooper of Culross. When he decided the case of MacCormick v. Lord Advocate as Lord President of the Court of Session, he stated, "the principle of unlimited sovereignty of Parliament is a distinctively English principle and has no counterpart in Scottish constitutional law." He continued, "Considering that the Union legislation extinguished the Parliaments of Scotland and England and replaced them by a new Parliament, I have difficulty in seeing why the new Parliament of Great Britain must inherit all the peculiar characteristics of the English Parliament but none of the Scottish." Nevertheless, he did not give a conclusive opinion on the subject. Thus, the question of Parliamentary sovereignty appears to remain unresolved. Parliament has not passed any Act defining its own sovereignty. Parliament's power has often been eroded by its own Acts. Acts passed in 1921 and 1925 grant the Church of Scotland complete independence in ecclesiastical matters. More recently, its power has been restricted by the United Kingdom's membership of the European Union, which has the power to make laws enforceable in each member state. In the Factortame case, the European Court of Justice ruled that UK courts could have powers to overturn legislation contravening EU law. This new power is a breach of parliamentary sovereignty, which is part of the UK constitution. Parliament has also created national devolved assemblies with legislative authority in Scotland, Wales and Northern Ireland. Similarly, it has granted the power to make regulations to Ministers of the Crown, and the power to enact religious legislation to the General Synod of the Church of England. (Measures of the General Synod and, in some cases, proposed statutory instruments made by ministers must be approved by both Houses before they become law.) In every case aforementioned, however, authority has been conceded by Act of Parliament, and may be taken back in the same manner. It is entirely within the authority of Parliament to, for example, abolish the devolved governments in Scotland, Wales and Northern Ireland or to leave the EU. However, especially in the case of withdrawing from EU membership, the political costs (the UK's economy and reputation in Europe would most likely be hugely damaged) of such a move would surely prevent it from occurring. Although legally Parliament's sovereignty has not been curtailed, in a politically sense it sovereignty has been reduced by Parliament's own Acts, especially the European Communities Act 1972 (UK), which made the UK a member of the EU. One well-recognised exception to Parliament's power involves binding future Parliaments. No Act of Parliament may be made secure from amendment or repeal by a future Parliament. For example, although the Act of Union 1800 states that the Kingdoms of Great Britain and Ireland are to be united "forever," Parliament permitted Southern Ireland to separate into a distinct nation, the Irish Free State, in 1922. Each House of Parliament possesses and guards various ancient privileges. The House of Lords rely on inherent right. In the case of the House of Commons, the Speaker goes to the Lords' Chamber at the beginning of each new Parliament and requests representatives of the Sovereign to confirm the Lower House's "undoubted" privileges and rights. The ceremony observed by the House of Commons dates to the reign of Henry VIII. Each House is the guardian of its privileges, and may punish breaches thereof. The extent of parliamentary privilege is based on law and custom. Sir William Blackstone states that these privileges are "very large and indefinite," and cannot be defined except by the Houses of Parliament themselves. The foremost privilege claimed by both Houses is that of freedom of speech in debate; nothing said in either House may be questioned in any court or other institution outside Parliament. Another privilege is that of freedom from arrest except for high treason, felony or breach of the peace; it applies from during a session of Parliament, as well as forty days before or after such a session. Members of both Houses are also privileged from service on juries. Both Houses possess the power to punish breaches of their privilege. Contempt of Parliament — for example, disobedience of a subpoena issued by a committee — may also be punished. The House of Lords may imprison an individual for any fixed period of time, but an individual imprisoned by the House of Commons is set free upon prorogation. The punishments imposed by either House may not be challenged in any court. - Blackstone, Sir William. (1765). Commentaries on the Laws of England. Oxford: Clarendon Press. - Davies, M. (2003). Companion to the Standing Orders and guide to the Proceedings of the House of Lords, 19th ed. (http://www.parliament.the-stationery-office.co.uk/pa/ld/ldcomp/compso.htm) - Farnborough, Thomas Erskine, 1st Baron. (1896). Constitutional History of England since the Accession of George the Third, 11th ed. London: Longmans, Green and Co. - "Parliament." (1911). Encyclopędia Britannica, 11th ed. London: Cambridge University Press. - The United Kingdom Parliament. Home Page. (http://www.parliament.uk/)
The Really Useful Book of Secondary Science Experiments 101 Essential Activities to Support Teaching and Learning How can a potato be a battery? How quickly will a shark find you? What food should you take with you when climbing a mountain? The Really Useful Book of Secondary Science Experiments presents 101 exciting, ‘real-world’ science experiments that can be confidently carried out by any KS3 science teacher in a secondary school classroom. It offers a mix of classic experiments together with fresh ideas for investigations designed to engage students, help them see the relevance of science in their own lives and develop a passion for carrying out practical investigations. Covering biology, chemistry and physics topics, each investigation is structured as a problem-solving activity, asking engaging questions such as, ‘How can fingerprints help solve a crime?’, or ‘Can we build our own volcano?’ Background science knowledge is given for each experiment, together with learning objectives, a list of materials needed, safety and technical considerations, detailed method, ideas for data collection, advice on how to adapt the investigations for different groups of students, useful questions to ask the students and suggestions for homework. Additionally, there are ten ideas for science based projects that can be carried out over a longer period of time, utilising skills and knowledge that students will develop as they carrying out the different science investigations in the book. The Really Useful Book of Secondary Science Experiments will be an essential source of support and inspiration for all those teaching in the secondary school classroom, running science clubs and for parents looking to challenge and excite their children at home. Table of Contents Experiment 1: Observation: Are probiotic yogurts worth the extra money? Experiment 2: Observation: How similar are animal and plant DNA? Experiment 3: Observation: What do the inside of lungs look like? Experiment 4: Observation: Are all fats the same? Experiment 5: Observation: How do plants exchange gases? Experiment 6: Observation: How do apples decay? Experiment 7: Is salt a good preserver of food? Experiment 8: Fair testing: How can plants use wind to reproduce? Experiment 9: Fair testing: Are there enzymes in our liver? Experiment 10: Fair testing: What is the best food to take with you when climbing a mountain? Experiment 11: Fair testing: Which is the most dangerous see to swim in if you are bleeding? Experiment 12: Fair testing: How quickly will our muscles tire? Experiment 13: Fair testing: Can we speed up the rate of photosynthesis? Experiment 14: Pattern seeking: Where do daises grow? Experiment 15: Pattern seeking: Do taller people have larger hands? Experiment 16: Pattern seeking: Do insects prefer to live in the light or the dark? Experiment 17: Pattern seeking: Can long legs jump further? Experiment 18: Pattern seeking: Do our hearts beat faster when we work harder? Experiment 19: Pattern seeking: Are hand-dryers more hygienic than paper towels? Experiment 20: Classification and identification: Can you identify animal and plant cells just by looking at them? Experiment 21: Classification and identification: Can we classify leaves? Experiment 22: Classification and identification: What’s the best fruit and vegetable to eat when you have a cold? Experiment 23: Classification and identification: How can fingerprints solve a crime? Experiment 24: Classification and identification: Which plants are growing near our school? Experiment 25: Classification and identification: What’s in our food? Experiment 26: Modelling: Can we build a digestive system? Experiment 27: Modelling: Can we build a DNA separating chamber? Experiment 28: Modelling: Can we build a model of DNA? Experiment 29: Modelling: Can we ferment our own ginger beer? Experiment 29: Modelling: Can we build a bug hotel? Experiment 30: Modelling: Can we design and make a stethoscope? Experiment 31: Observation: Can a solid turn into a gas? Experiment 32: Observation: Where should we dig for oil?Experiment 33: Observation: What colour are M&Ms? Experiment 34: Observation: What is the best material for a campfire? Experiment 35: Observation: How can we make colourful flames? Experiment 36: Observation: What is special about the melting and freezing point of a substance? Experiment 37: Fair Testing: Which is the best washing powder? Experiment 38: Fair Testing: Can we prevent rusting? Experiment 39: Fair Testing: Which antacid is the most effective? Experiment 40: Fair Testing: Which is the best brand of disposable nappies? Experiment 41: Fair Testing: How does temperature affect the rate of a reaction? Experiment 42: Fair Testing: How quickly will a puddle evaporate on a hot day? Experiment 43: Pattern seeking: How quickly will a battery run down? Experiment 44: Pattern Seeking: What is the hardest liquid to swim through? Experiment 45: Pattern Seeking: Will aquatic plants grow in acidic water? Experiment 46: Pattern Seeking: Do all oxides have the same pH? Experiment 47: Pattern Seeking: Which element in group 2 of the periodic table is the most reactive? Experiment 48: Pattern Seeking: Which element in group 7 of the periodic table is the most reactive? Experiment 49: Classification and Identification: Are all changes reversible? Experiment 50: Classification and Identification: What is the best soil for growing plants? Experiment 51: Classification and Identification: How can we identify colourless gases? Experiment 52: Classification and Identification: How can polymers be identified? Experiment 53: Classification and Identification: Do chemical reactions always give off heat? Experiment 54: Classification and Identification: Does everything dissolve in water? Experiment 55: Modelling: Can we make our own fizzing bath bombs? Experiment 56: Modelling: Can we make popping fruit juice balls? Experiment 57: Modelling: Can we grow a crystal garden? Experiment 58: Modelling: Can we build our own volcano? Experiment 59: Modelling: How can cabbage be an indicator? Experiment 60: Modelling: Can we make a bouncing custard ball? Experiment 61: Observation: How many colours are there in light? Experiment 62: Observation: How does pressure vary in a water column? Experiment 63: Observation: What do waves look like? Experiment 64: Observation: Which objects will give you a static shock? Experiment 65: Observation: How do gases move? Experiment 66: Observation: How much ‘stuff’ do we make in a reaction? Experiment 67: Fair Testing: How can we change the brightness of a bulb? Experiment 68: Fair Testing: Why do moon craters vary in size? Experiment 69: Fair Testing: What are the most dangerous weather conditions to drive in? Experiment 70: Fair Testing: How can we increase the resistance in a circuit? Experiment 71: Fair Testing: How can blood spatter solve a crime? Experiment 72: Fair Testing: Can we stop radio waves? Experiment 73: Pattern seeking: How can you make a swing go faster? Experiment 74: Pattern Seeking: Can you break a spring? Experiment 75: Pattern Seeking: How can we make a magnet stronger? Experiment 76: Pattern Seeking: How does light enter and leave a mirror? Experiment 77: Pattern Seeking: How can we change the speed of light? Experiment 78: Pattern Seeking: What happens to waves in shallow water? Experiment 79: Classification and Identification: Which materials are best for keeping something warm? Experiment 80: Classification and Identification: Which materials are best for building an electric circuit? Experiment 81: Classification and Identification: Can we identity different types of radiation? Experiment 82: Classification and Identification: Can we classify all materials as solids, liquids or gases? Experiment 83: Classification and Identification: What is the densest liquid?Experiment 84: Classification and Identification: Where is the energy going? Experiment 85: Modelling: Can we cook food using the sun? Experiment 86: Modelling: Can we make our own camera? Experiment 87: Modelling: How can a potato be a battery? Experiment 88: Modelling: Can we build a catapult? Experiment 89: Modelling: Can we design and make a musical instrument? Experiment 90: Modelling: Can we make a crash helmet? Project 1: Healthy teeth Project 3: Environmental survey Project 4: Set Design Project 5: Olympic science Project 6: Chocolate lab Project 7: Scene of crime investigation Project 8: Fairground games Project 9: Aeroplane design Project 10: What’s the weather like? Tracy-ann Aston is Lecturer in Education and Teacher Training, specialising in Science Education and Primary Teacher Training at the University of Bedfordshire, UK "There is a perennial debate about how best to assess practical work in school science, with many repercussions for examinations and classroom practice... This book sets out to do just that for 101 experiments relating to ages 11–14 science."- Trevor Critchley, Education in Chemistry - Editable Tables
We know that astronomical data give us accurate values of the radii of the Sun, the Earth and the Moon. Furthermore, the knowledge of their relative distances predicts quite accurately the instant when the umbra-penumbra limit sweeps some specific craters on the Moon during lunar eclipses. Since the 1830s, crater timing has been used during lunar eclipses to measure the length of the Earth’s shadow. The method is simple: one takes the timing of lunar features (craters, limbs, ridges, peaks, bright spots) as they enter and exit the umbra. The Sun-Earth-Moon geometry being known quite precisely is then possible to calculate the size and shape of the Earth’s umbra at the Moon. Measurements that vary from one eclipse to the next can now be made with low-power telescopes or a clock synchronized with radio time signals. However, it has systematically been found that the shadow of the Earth seems to be 2% larger than what is expected from geometrical predictions. Even, if it is believed that the thickness of the Earth atmosphere is responsible for that displacement , it was realized that the atmospheric absorption cannot explain the absorption of light at a height of up to 90 km above the Earth, as required by this hypothesis. It may be noted in particular that Link has firmly established a relationship between the enlargement of the Earth’s shadow during lunar eclipses and the presence of meteors, which have the ability to distort the optical properties of the atmosphere when they are braked at high altitudes . It has been said that the pronounced red colour in the inner portions of the umbra during an eclipse of the Moon is caused by refraction of sunlight through the upper regions of the Earth’s atmosphere, but the umbral shadow towards the centre is too bright to be accounted for by refraction of visible sunlight. In Sections 2 and 3, we give a brief history of the enlargement of the Earth’s umbra and the excess of light into the Earth’s shadow onto the Moon during lunar eclipses. We present some accepted interpretations and we show how the Allais effect, which occurs at the time when problems arise related to these anomalies, leads us to reject these interpretations. In Section 4, it emerges from a discussion that, failing to have an answer that would explain the two coexisting anomalies, the Allais eclipse effect currently remains the only viable option. Experiments are proposed as much to corroborate the observations of the two anomalies as to test the Allais eclipse effect. We conclude that both anomalies during lunar eclipses are caused by a lunar Allais effect. 2. Enlargement of the Earth’s Umbra 2.1. Brief History of the Enlargement of the Earth’s Umbra on the Moon during Lunar Eclipses In the early 1700s, Philippe de La Hire made a curious observation about Earth’s umbra. The predicted radius of the shadow needed to be enlarged by about 1/41 in order to fit timings made during an eclipse of the Moon (La Hire 1707). Additional observations over the next two centuries revealed that the shadow enlargement was somewhat variable from one eclipse to the next . Chauvenet (1891) adopted a value of 1/50, which has become the standard enlargement factor for lunar eclipse predictions published by many national institutes worldwide. Some authorities dispute Chauvenet’s shadow enlargement convention . Danjon (1951) notes that the only reasonable way of accounting for a layer of opaque air surrounding Earth is to increase the planet’s radius by the altitude of the layer . This can be accomplished by proportionally increasing the parallax of the Moon. The radii of the umbral and penumbral shadows are then subject to the same absolute correction and not the same relative correction employed in the traditional Chauvenet 1/50 convention. Danjon estimates the thickness of the occulting layer to be 75 km and this results in an enlargement of Earth’s radius and the Moon’s parallax of about 1/85. Since 1951, the French almanac Connaissance des Temps has adopted Danjon’s method for the enlargement Earth’s shadows in their eclipse predictions Danjon’s method correctly models the geometric relationship between an enlargement of Earth’s radius and the corresponding increase in the size of its shadows. Meeus and Mucke (1979), and Espenak (2006), both use Danjons method. However, the resulting umbral and penumbral eclipse magnitudes are smaller by approximately 0.006 and 0.026 respectively as compared to predictions using the traditional Chauvenet convention of 1/50. For his part, in an analysis of 57 eclipses over a period of 150 years, Link (1969) found an enlargement of the shadow of 2.3% on average. Furthermore, schedules inputs and outputs of the crater through the umbra for four lunar eclipses from 1972 to 1982 strongly support the Chauvenet value of 2%. Of course, the small magnitude difference between the two methods is difficult to observe because the edge of the umbral shadow is diffuse. From a physical point of view, there is no well defined border between the umbra and the penumbra. The shadow density actually varies continuously as a function of radial distance from the central axis out to the extreme limit of the penumbral shadow. However, the density variation is most rapid near the theoretical edge of the umbra. Kuhl’s (1928) contrast theory demonstrates that the verge of the umbra is perceived at the point of inflexion in the shadow density. This point appears to be equivalent to a layer in Earth’s atmosphere at an altitude of about 120 to 150 km. The net enlargement of Earth’s radius of 1.9% to 2.4% corresponds to an extension of the umbra of 1.5%, to 1.9%, in reasonably good agreement with the conventional value. It seems that the increase of the Earth’s umbral shadow during eclipses of the Moon is the classical value of 2% (the rule of the fiftieth) used in most calculations of lunar eclipses . 2.2. Accepted Interpretation of the Enlargement of the Umbra Numerous reports show that the umbra-penumbra limit appears significantly displaced on the moon during an eclipse. It is believed that the thickness of the Earth atmosphere is responsible for that displacement . In order to study more deeply the phenomenon showing that the umbra-penumbra limit appears significantly displaced on the Moon during an eclipse, it is important to evaluate if the reported increase of 2% of the Earth’s shadow on the Moon corresponds to a reasonable value of the height at which the atmosphere is opaque. Calculations indicate that this enlargement corresponds to a terrestrial altitude of 92 km. This usual interpretation of the umbral enlargement forces us to believe that the atmosphere is normally opaque up to 92 km or so. But how is this possible when, at this altitude, the air is extremely rarefied? It is the altitude close to the orbit on which a satellite travels around the Earth. In fact, according to data , the atmospheric pressure at 90 km above sea level is about half a million times smaller than that at sea level. Above 15 km, the atmosphere becomes relatively transparent to light, since 90% of the air and almost all the humidity and pollution are below that level. That makes an enlarged obscuration due to the opacity of the atmosphere of only 0.3% which is much smaller than the 2.0% reported. Furthermore, the eruption of volcanos cannot explain the larger shadow. According to some, the altitude reached by some material ejected from volcano El Chichon is in the stratosphere, some 26 kilometers (16 miles) above Earth’s surface – roughly 50% higher than material from the famous Mount St. Helens . Since the atmosphere does not appear to be responsible for the umbra-penumbra limit displacement of 2% on the Moon, then what is the cause? F. Link argues that the meteoric dust in the upper atmosphere of the Earth is at the origin of the additional weakening of the light and the expansion of the Earth’s darkness . We might point out, in particular, that Link has actually established a concomitance between the enlargement of the Earth’s umbra during lunar eclipses and the presence of meteors, which are capable of distorting the optical properties of the atmosphere when they are decelerated at high elevations . Paul Marmet and Christine Couture , for their part, believe that the actual umbra of the Earth projected on the Moon is not as big as observed, that the sensitivity of the eyes is a factor leading necessarily to an umbral enlargement and that almost the totality of the reported umbra-penumbra limit displacement is an optical effect that has nothing to do with the thickness of the Earth atmosphere. For our part, we believe that the observed times to browse the path of the Moon through the Earth’s obscurity deviate from the predicted times and that some variations in colour, size and shape of the umbra occur in the darkness. We attribute this deviations and variations to the Allais eclipse effect. 2.3. Umbral Enlargement and the Allais Eclipse Effect During an eclipse of the Moon, it is predicted geometrically that the photons from the Sun describe a rectilinear trajectory as if they were little deflected and pass at a minimum approach distance from the centre of the Earth (slightly larger than the radius of the Earth), before moving to the Moon. A ray of sunlight passes close to the Moon at point of minimum approach to arrive at the point P at the end of the shadow cone of the Earth. The trajectory followed by the solar photons shapes the curvature of minimum approach of the Earth and the Moon. However, we have serious reasons to believe that during a lunar eclipse, with the Earth interposed between the Moon and the Sun, there would be a kind of anti-gravity on the Moon which would be manifested by a deviation of light. This is precisely the Allais effect . The orbital radius seems longer, which means that the curvature of minimum approach of the Earth and the Moon is shifted outward. Since the curvature is the inverse of the square of the radius, the curvature is even smaller than the radius is large. This is grounded in the Newtonian logic stating that gravity identified to the curvature is all the more weak as the orbital radius is large. The deviated photons will pass at a distance of minimum approach from the centre of the Earth and at the point of minimum approach of the Moon to end at the point , casting an enlarged umbra cone. It matches with the observed cone of the enlargement of the Earth’s umbra. During a lunar eclipse, it is predicted geometrically that the photons from the Sun describe a rectilinear trajectory as if they were a little deflected (Figure 1). Inevitably, sunlight, observed on the eclipsed Moon, will tend to move away. The photon has an “inertial mass” equal to equivalent to a “gravitational mass” also equal to . The energy of the “fallen” photon from the Sun will be instead of mgH (m = inertial and gravitational mass of the photon; g = acceleration due to gravity, H = height). The gravitational mass of photons, lessened, takes distance, so increasing the darkness. We are witnessing an abnormal gravitational frequency shift. Suppose that in normal times the light is emitted by the Sun at the height H (distance Sun-moon) . The total energy of a photon of frequency v and Figure 1. Exaggerated diagram (for comprehension) of the cones of the umbra. The illumination of the Earth by the Sun projects into space a converging cone of umbra and a divergent cone of penumbra whose conical generators are tangent to the Earth and the Sun. We represent here that the converging cone of umbra to illustrate the enlargement of the earth’s shadow. The umbra cone does not completely obscure the Moon and, as early as the 18th century, astronomers knew that the shadow limit was a little further ( ) than according to the geometric path of the rays (P). The quasi-parallel interior tangents to the Sun and the Earth give the two interior cones, predicted and observed, of the umbra. The tangent gives the observed cone of the umbra while the tangent gives the calculated cone of the umbra. energy hv, reaching the lunar surface has become The receiver, placed on the lunar ground, detects a frequency greater than v of the solar source (g designates the lunar gravitational field): During a lunar eclipse, due to a potential loss of attraction, the lunar gravitational field g amounts to . When a photon emitted by the Sun reaches the surface of the Moon, he lost potential energy and won the kinetic energy . Its total energy has become The frequency of the photon at its arrival at the surface of the eclipsed Moon is less red-shifted relative to its initial frequency, according to the relation The receiver detects a frequency slightly smaller than of not eclipsed Moon. This means a small blue shift for the Sun during the eclipse. If this hypothesis is correct which consists to declare that the Allais effect causes a kind of repulsion between the three celestial bodies involved, there should be a variation of the gravitational potential. This means that the gravitation will influence the geometry of space-time: the time of clocks and the length measured by a rule will be affected depending on whether there is more or less gravity. Einstein’s general theory of relativity predicts that a clock in the presence of weak gravity runs more rapidly than one located where gravity is stronger. Consequently, the frequencies of radiation emitted by atoms in the presence of a weak gravitational field are shifted to higher frequencies when compared with the same emissions in the presence of a strong field. The light of the Sun observed on an eclipsed Moon should be blue shifted; a fraction of the solar gravitational redshift ( ) which is about two parts in a million . The widening of the Earth’s shadow on the Moon would not be due to a greater density of the upper atmosphere of the Earth, which would make it as opaque as the lower atmosphere, it would be caused by a gravitational potential temporarily decreased. An “alleviated” matter would dictate to space-time a smaller degree of curvature; the space-time would in turn impose to matter to move on a larger orbital radius. 3. Excess of Light into the Earth’s Shadow 3.1. Brief History of the Excess of Light into the Earth’s Shadow on the Moon during Lunar Eclipses The first work on the variation in brightness of eclipses was executed by André-Louis Danjon in 1920. He devised a scale of brightness for total lunar eclipses, from 0 for invisible to 4 for very brilliant. He used it to analyze data on eclipses extending back over three and a half centuries, and showed a correlation between the eclipse brightness and the solar activity. But a series of three-color photometric observations of Moon eclipse, made by him and his associates between 1932 and 1957, appears to show a clear correlation between the eclipse brightness and the geomagnetic planetary index Kp . So the new data seems to contradict the first Danjon’s conclusion. An attempt to interpret the relation with Kp in terms of lunar luminescence indicates that the change in the eclipse brightness is in accord with the rate of increase of the plasma energy, as predicted by the experiment of Snyder et al. (1963) . However, there is some difficulty in the required proton density being greater than the observed value by an order of magnitude. The same calculations show that luminescence to be visible in ordinary moonlight requires a plasma energy at least three orders of magnitude greater than the maximum value predicted by Snyder et al. They also show that the reported dates of these observations fall on geophysically quiet days, as well as on dates of high Kp. The above conclusion agrees with the result of calculations by Ney et al. in 1966. On the other hand, J. Dubois and F. Link in 1969 found a correlation between the brightness of the eclipsed Moon and the solar activity, as had been suggested by Danjon on the basis of its first observations. They demonstrated that the brightness of eclipse was related not only to the sunspots number but also to the height of the latitude. They showed an annual correlation between the heliographic latitude of the apparent centre of the Sun’s disk and eclipse brightness . It was suggested that the brightness anomaly of the umbral region during an eclipse of the Moon would be caused by refraction of sunlight through the upper regions of the Earth’s atmosphere. The red coloring arises because, they say, sunlight reaching the Moon must pass through a long and dense layer of the Earth’s atmosphere, where it is scattered. Shorter wavelengths are more likely to be scattered by the small particles and so, by the time the light has passed through the atmosphere, the longer wavelengths dominate. This resulting light we perceive as red. The amount of refracted light depends on the amount of dust or clouds in the atmosphere; this also controls how much light is scattered. In general, the dustier the atmosphere, the more that other wavelengths of light will be removed (compared to red light), leaving the resulting light a deeper red color . Despite this reasoning, it has been found that towards the centre the umbra is too bright to be accounted for by refraction of visible sunlight. F. Link proposed that this excess be interpreted as luminescence . He concluded that about 10 percent of the Moon’s optical radiation is caused by luminescence. Observations seem to confirm the existence of lunar luminescence. The term luminescence can be applied to any object that emits light in addition to the usual reflected light . The main characteristic of luminescence is that the emitted light is an attribute of the object itself, and the light emission is stimulated by some internal or external process. As external process, Link suggested the luminescence of the lunar surface by X-ray bombardment from the uneclipsed regions of the solar corona, as suggested by Link. This theory is supported by the variation of a factor of 100, between solar maximum and minimum, of the intensity of certain wavelengths of X-rays . Another possible mechanism by which the eclipsed Moon shines results from the fact that the Moon is covered by a fine layer of meteoric dust, and would therefore contain quantities of the achondritic enstatites. This type of stony meteorite produces the luminescence when protons and electrons of the solar wind are deflected and impinge on the lunar surface during a total eclipse. An experiment performed by Zdenek Kopal, C. J. Derham and J. E. Geakel in 1963 showed that certain meteorite specimens glowed with a strange red light, same colour as the umbra in eclipse, when bombarded by high energy protons in the laboratory . It appears to us that the excess of irradiation of the Moon in the shadow of the Earth during the eclipse is partially caused by refraction in the atmosphere, but that it prevailingly depends of the light emission stimulated by an internal process linked to the Allais eclipse effect. 3.2. Luminescence of the Eclipsed Moon and Allais Effect The time of vibration (T) of atoms and molecules of luminescent gases in the fieldless region of space is (T is the time of vibration in the atom at rest; is the modified time of vibration) . In a region of space with the gravitational field, the time of vibration is altered to [ contains the gravitational potential ( )]. Our assumption is that at the surface of an eclipsed Moon there is a weaker gravitational field than in a frame of reference without eclipse ( ) During the lunar eclipse time, the gravitational potential of the Moon is conjecturally diminished; the time of the vibration of the atom is shorter. The metric of the obscured Moon is affected and the ticking of time accelerates relative to the system without eclipse in which the atom is considered at rest. Consequently, the red shift of the spectral lines of light that comes from the layer of particles on the ground will have a small additional blue “Allais” shift which reduces the “Einstein” redshift . As the Einstein effect (i.e. the tiny frequency shift of spectral lines in a gravitational field) is directed towards the blue, there is thus more internal electromagnetic energy. It appears that this blueshift by variation in the reduction of the mass could be a form of atomic excitement at the level of electrons, as Brownian motion. More specifically, we would say that the Allais eclipse effect would have engendered a significant change of wavelength within the molecules of matter . The excess of luminescence would be the imprint left on the light by the intramolecular oscillation of the atoms constituting the molecules of the lunar soil which spreads it. A Raman effect caused by an Allais effect, in some way. 3.3. An Anti-Stoke Raman Effect Induced by an Allais Effect We assume that the lunar gravitational potential can be reduced in times of eclipse, what would accelerate the vibration of atoms. A molecule can be excited to a very high energy state. The amount of energy necessary to reach this excited state is . Therefore, the relaxation of the molecule to the ground-state vibrational energy level results in the emission of a photon of energy . This emission is usually observed in the visible spectral region and is called Rayleigh scattering. We think that the rapid oscillation of a light wave passing by the intramolecular level of atoms which constitute the molecules diffused by the lunar soil could be similar to an effect Raman anti-Stoke . The scattering of light on the optical modes is designated Raman effect. It is different from the Rayleigh scattering because the scattered light changes the frequency of the spectrum active vibration. Historically, the effect was first observed with molecules. Molecules vibrate, and each molecular oscillation corresponds to a certain amount of energy. In the scattering process, this energy is added or subtracted from the incident light. An anti-Stokes Raman effect occurs when the molecule absorbs an incident light of frequency and reemits it at a higher frequency. Thus, during the eclipse of the Moon, the excited molecule would oscillate from a superior vibrational energy level, say . The energy absorbed in this process is still . The molecule can relax to the original vibration energy level and emits a photon ; however, the relaxation can be to the ground state. The return to the state results in the emission of a photon which is greater than the exciting energy from level 1. The photon energy emitted is . Spectral lines with frequencies higher than are labeled anti-Stokes lines . This Raman shift induces a brighter electromagnetic radiation. 4. Discussion and Conclusion In 2009, NASA’s Lunar Reconnaissance Orbiter (LRO) was launched with the Lunar Crater Observation and Sensing Satellite (LCROSS) on the first U.S. mission to the Moon in over 10 years. LRO gathered information on day-night temperature maps, contributed data for a global geodetic grid, and conducted high-resolution imaging. During lunar eclipses, the solar-powered orbiter also falls in Earth’s shadow, cutting it off from the source of its power. The mission controllers can then use an instrument—called Diviner—that can watch how the lunar surface responds to the rapid change in temperature caused by a lunar eclipse. These data which help scientists better understand the composition and properties of the surface could be a scientific boon for understanding both anomalies . We believe that almost the totality of the reported umbra-penumbra boundary shift and the excessive clarity of the penumbra reveal a lunar Allais effect on the Moon’s shadow that has nothing to do with the thickness of the Earth’s atmosphere. Both phenomena were reported during each and every lunar eclipse recorded for the past 180 years. They occur during lunar eclipses and are correlated. Dr Marmet demonstrated that the Earth’s atmosphere cannot be the cause of the enlargement of the Earth’s shadow. He concludes that it is an optical illusion, but neither addresses nor explains the second offset: the excessive brightness of the penumbra . If he is right to say that the Earth’s atmosphere is not responsible for the 2% umbra-penumbra limit shift on the Moon, he is wrong to evoke the optical illusion. NASA records anomalies without providing an explanation, the priorities being elsewhere. The door is open to researchers to probe the reasons and suggest fields to explore. The list of our references shows that they cannot be explained by current science, which leaves only one option: the lunar Allais effect. But the bottleneck, which means that this aspect of science remains speculative even as Professor Allais’ experiments have validated the solar eclipse effect, is the question of a lack of willingness to experiment. How can interest be aroused in experimenters for whom the scientific value of precise experience is dependent on their theoretical interpretation? Classical conservative physical thought cannot tolerate the defeat of the current theory of gravitation when applied to the case of the influence of the attraction of the Sun and the Moon on the motion of the paraconic pendulum, whether these are the amplitudes of the lunisolar periodic components or the anomalies observed during the total eclipses of the Sun . However, the most accommodating physicists say they do not rely on the more or less contradictory experiments carried out so far; they would like experiments operated with a paraconic pendulum at any point similar to the pendulum used by M. Allais, or they would like to turn to more radical experiments, like those intended for modern theories. For example, the atomic clock cooled by cesium laser (PHARAO) placed by the European Space Agency (ESA) on the International Space Station (ISS) could have been used. In default of confirm doubtful fashionable theories, this high technology could test the Allais effect and supply the way to tie, by a new theoretical link, the facts observed during the eclipses to the physical laws having received the sanction of a rigorous experimental control. In conclusion, it seems that two noticed anomalies during lunar eclipses, the enlargement of the Earth’s shadow delineated onto its satellite and an excessive illumination of the penumbra, adjusted ad hoc to the Earth’s atmosphere, would rather be caused by an Allais eclipse effect, i.e., a repulsion that occurs when the Moon passes directly behind the Earth into its umbra, when Sun, Earth and Moon are closely aligned in space. This is consistent with our knowledge of the solar eclipse, with the calculation of the abnormal spontaneous acceleration of the Moon during the solar eclipse in June 1954 (paraconical pendulum of Maurice Allais) and the result recorded by a gravimeter during the solar eclipse of 1997. (To know more about lunar eclipses by pictures, references have been added). Link, F. and Linková, Z. (1954) Agrandissement de l’ombre terrestre pendant les éclipses de Lune; influences météoriques. Publishing House of the Czechoslovak Academy of Sciences. Provided by the NASA Astrophysics Data System. Dubois, J. and Link, F. (1970) Analyse photométrique de l’ombre intérieure pendant les éclipses de Lune. Publishing House of the Czechoslovak Academy of Sciences. Provided by the NASA Astrophysics Data System. Sandulak, N. (1964) Jurgen Stock, Indication of Luminescence Found in the December 1964 Lunar Eclipse. Cerro Tololo Inter-American Observatory, La Serena, 237. (1965) Provided by the NASA Astrophysics Data System.
Ubiquitous on Earth, water also has been found in comets, on Mars and in molecular clouds in interstellar space. Now, scientists say this common fluid is not as well understood as we thought. “Water, as we know it, does not exist within our bodies,” said Martin Gruebele, a William H. and Janet Lycan Professor of Chemistry at the University of Illinois. “Water in our bodies has different physical properties from ordinary bulk water, because of the presence of proteins and other biomolecules. Proteins change the properties of water to perform particular tasks in different parts of our cells.” Consisting of two hydrogen atoms and one oxygen atom, water molecules are by far the body’s largest component, constituting about 75 percent of body volume. When bound to proteins, water molecules participate in a carefully choreographed ballet that permits the proteins to fold into their functional, native states. This delicate dance is essential to life. “While it is well known that water plays an important role in the folding process, we usually only look at the motion of the protein,” said Gruebele, who also is the director of the U. of I.’s Center for Biophysics and Computational Biology, and a researcher at the Beckman Institute. “This is the first time we’ve been able to look at the motion of water molecules during the folding process.”Using a technique called terahertz absorption spectroscopy, Gruebele and his collaborator Martina Havenith at the Ruhr-University Bochum studied the motions of a protein on a picosecond time scale (a picosecond is 1 trillionth of a second). The researchers present their findings in a paper published July 23 in the online version of the chemistry journal Angewandte Chemie. Terahertz spectroscopy provides a window on protein-water rearrangements during the folding process, such as breaking protein-water-hydrogen bonds and replacing them with protein-protein-hydrogen bonds, Gruebele said. The remaking of hydrogen bonds helps organize the structure of a protein. In tests on ubiquitin, a common protein in cells, the researchers found that water molecules bound to the protein changed to a native-type arrangement much faster than the protein. The water motion helped establish the correct configuration, making it much easier for the protein to fold. “Water can be viewed as a ‘designer fluid’ in living cells,” Gruebele said. “Our experiments showed that the volume of active water was about the same size as that of the protein.” The diameter of a single water molecule is about 3 angstroms (an angstrom is about one hundred-millionth of a centimeter), while that of a typical protein is about 30 angstroms. Although the average protein has only 10 times the diameter of a water molecule, it has 1,000 times the volume. Larger proteins can have hundreds of thousands times the volume. A single protein can therefore affect, and be influenced by, thousands of water molecules.“We previously thought proteins would affect only those water molecules directly stuck to them,” Gruebele said. “Now we know proteins will affect a volume of water comparable to their own. That’s pretty amazing.” Funding was provided by the Human Frontier Science Program and the National Science Foundation. James E. Kloeppel \| University of Illinois Flow of cerebrospinal fluid regulates neural stem cell division 22.05.2018 \| Helmholtz Zentrum München - Deutsches Forschungszentrum für Gesundheit und Umwelt Chemists at FAU successfully demonstrate imine hydrogenation with inexpensive main group metal 22.05.2018 \| Friedrich-Alexander-Universität Erlangen-Nürnberg So-called quantum many-body scars allow quantum systems to stay out of equilibrium much longer, explaining experiment \| Study published in Nature Physics Recently, researchers from Harvard and MIT succeeded in trapping a record 53 atoms and individually controlling their quantum state, realizing what is called a... The historic first detection of gravitational waves from colliding black holes far outside our galaxy opened a new window to understanding the universe. A... A team led by Austrian experimental physicist Rainer Blatt has succeeded in characterizing the quantum entanglement of two spatially separated atoms by observing their light emission. This fundamental demonstration could lead to the development of highly sensitive optical gradiometers for the precise measurement of the gravitational field or the earth's magnetic field. The age of quantum technology has long been heralded. Decades of research into the quantum world have led to the development of methods that make it possible... Cardiovascular tissue engineering aims to treat heart disease with prostheses that grow and regenerate. Now, researchers from the University of Zurich, the Technical University Eindhoven and the Charité Berlin have successfully implanted regenerative heart valves, designed with the aid of computer simulations, into sheep for the first time. Producing living tissue or organs based on human cells is one of the main research fields in regenerative medicine. Tissue engineering, which involves growing... A team of scientists of the Max Planck Institute for the Structure and Dynamics of Matter (MPSD) at the Center for Free-Electron Laser Science in Hamburg investigated optically-induced superconductivity in the alkali-doped fulleride K3C60under high external pressures. This study allowed, on one hand, to uniquely assess the nature of the transient state as a superconducting phase. In addition, it unveiled the possibility to induce superconductivity in K3C60 at temperatures far above the -170 degrees Celsius hypothesized previously, and rather all the way to room temperature. The paper by Cantaluppi et al has been published in Nature Physics. Unlike ordinary metals, superconductors have the unique capability of transporting electrical currents without any loss. Nowadays, their technological... 02.05.2018 \| Event News 13.04.2018 \| Event News 12.04.2018 \| Event News 18.05.2018 \| Power and Electrical Engineering 18.05.2018 \| Information Technology 18.05.2018 \| Information Technology
Economic history of Germany Part of a series on the \|History of Germany\| Germany before 1800 was heavily rural, with some urban trade centers. In the 19th century it began a stage of rapid economic growth and modernization, led by heavy industry. By 1900 it had the largest economy in Europe, a factor that played a major role in its entry into World War I and World War II. Devastated by World War II, West Germany became an "economic miracle" in the 1950s and 1960s with the help of the Marshall Plan. Currently it is the largest individual economy in the EU with GDP of roughly 3 trillion USD. - 1 Middle Ages - 2 Early modern era - 3 Industrial revolution - 4 20th century - 5 Nazi economy - 6 Post-World War II - 7 Social market economy - 8 Economic miracle and beyond - 9 German reunification and its aftermath - 10 See also - 11 Notes - 12 Further reading Medieval Germany, lying on the open Northern European Plain, was divided into hundreds of contending kingdoms, principalities, dukedoms, bishoprics, and free cities. Economic prosperity did not mean geographical expansion; it required collaboration with some, competition with others, and an intimate understanding among government, commerce, and production. A desire to save was also born in the German experience of political, military, and economic uncertainty. Towns and cities The German lands had a population of about 5 or 6 million. The great majority were farmers, typically in a state of serfdom under the control of nobles and monasteries. A few towns were starting to emerge. From 1100, new towns were founded around imperial strongholds, castles, bishops' palaces and monasteries. The towns began to establish municipal rights and liberties (see German town law). Several cities such as Cologne became Imperial Free Cities, which did not depend on princes or bishops, but were immediately subject to the Emperor. The towns were ruled by patricians (merchants carrying on long-distance trade). The craftsmen formed guilds, governed by strict rules, which sought to obtain control of the towns; a few were open to women. Society was divided into sharply demarcated classes: the clergy, physicians, merchants, various guilds of artisans; full citizenship was not available to paupers. Political tensions arose from issues of taxation, public spending, regulation of business, and market supervision, as well as the limits of corporate autonomy. Cologne's central location on the Rhine river placed it at the intersection of the major trade routes between east and west and was the basis of Cologne's growth. The economic structures of medieval and early modern Cologne were characterized by the city's status as a major harbor and transport hub upon the Rhine. It was governed by its burghers. It was a business alliance of trading cities and their guilds that dominated trade along the coast of Northern Europe and flourished from the 1200 to 1500, and continued with lesser importance after that. The chief cities were Cologne on the Rhine River, Hamburg and Bremen on the North Sea, and Lübeck on the Baltic. The Hanseatic cities each had its own legal system and a degree of political autonomy. Early modern era Thirty Years War The Thirty Years' War (1618–1648) was ruinous to the twenty million civilians and set back the economy for generations, as marauding armies burned and destroyed what they could not seize. The fighting often was out of control, with marauding bands of hundreds or thousands of starving soldiers spreading plague, plunder, and murder. The armies that were under control moved back and forth across the countryside year after year, levying heavy taxes on cities, and seizing the animals and food stocks of the peasants without payment. The enormous social disruption over three decades caused a dramatic decline in population because of killings, disease, crop failures, declining birth rates and random destruction, and the out-migration of terrified people. One estimate shows a 38% drop from 16 million people in 1618 to 10 million by 1650, while another shows "only" a 20% drop from 20 million to 16 million. The Altmark and Württemberg regions were especially hard hit. It took generations for Germany to fully recover. According to John Gagliardo, the recovery period lasted for about fifty years until the end of the century and was over by the 1700s. At that time, Germany probably had reached its pre-war population (though this is disputed). Then, there was a period of steady though quite slow growth to the 1740s. Afterward came a period of rapid but not exceptional economic expansion, that mainly occurred in the great states in the east (Austria, Saxony, Prussia) rather than in the small states of central or south Germany. Peasants and rural life Peasants continued to center their lives in the village, where they were members of a corporate body and help manage the community resources and monitor the community life. Across Germany and especially in the east, they were serfs who were bound permanently to parcels of land. In most of Germany, farming was handled by tenant farmers who paid rents and obligatory services to the landlord, who was typically a nobleman. Peasant leaders supervised the fields and ditches and grazing rights, maintained public order and morals, and supported a village court which handled minor offenses. Inside the family the patriarch made all the decisions, and tried to arrange advantageous marriages for his children. Much of the villages' communal life centered around church services and holy days. In Prussia, the peasants drew lots to choose conscripts required by the army. The noblemen handled external relationships and politics for the villages under their control, and were not typically involved in daily activities or decisions. The emancipation of the serfs came in 1770-1830, beginning with then Danish Schleswig in 1780. Prussia abolished serfdom with the October Edict of 1807, which upgraded the personal legal status of the peasantry and gave them the chance to purchase for cash part of the lands they were working. They could also sell the land they already owned. The edict applied to all peasants whose holdings were above a certain size, and included both Crown lands and noble estates. The peasants were freed from the obligation of personal services to the lord and annual dues. A bank was set up so that landowner could borrow government money to buy land from peasants (the peasants were not allowed to use it to borrow money to buy land until 1850). The result was that the large landowners obtained larger estates, and many peasant became landless tenants, or moved to the cities or to America. The other German states imitated Prussia after 1815. In sharp contrast to the violence that characterized land reform in the French Revolution, Germany handled it peacefully. In Schleswig the peasants, who had been influenced by the Enlightenment, played an active role; elsewhere they were largely passive. Indeed, for most peasants, customs and traditions continued largely unchanged, including the old habits of deference to the nobles whose legal authority remained quite strong over the villagers. Although the peasants were no longer tied to the same land like serfs had been, the old paternalistic relationship in East Prussia lasted into the 20th century. Before 1850 Germany lagged behind the leaders in industrial development, Britain, France and Belgium. However, the country had considerable assets : a highly skilled labor force, a good educational system, a strong work ethic, good standards of living and a sound protectionist strategy based on the Zollverein. By midcentury, the German states were catching up, and by 1900 Germany was a world leader in industrialization, along with Britain and the United States. In 1800, Germany's social structure was poorly suited to any kind of social or industrial development. Domination by modernizing France during the era of the French Revolution (1790s to 1815) produced important institutional reforms, including the abolition of feudal restrictions on the sale of large landed estates, the reduction of the power of the guilds in the cities, and the introduction of a new, more efficient commercial law. Nevertheless, traditionalism remained strong in most of Germany. Until midcentury, the guilds, the landed aristocracy, the churches, and the government bureaucracies had so many rules and restrictions that entrepreneurship was held in low esteem, and given little opportunity to develop. From the 1830s and 1840s, Prussia, Saxony, and other states reorganized agriculture, introducing sugar beets, turnips, and potatoes, yielding a higher level of food production that enabled a surplus rural population to move to industrial areas. The beginning of the industrial revolution in Germany came in the textile industry, and was facilitated by eliminating tariff barriers through the Zollverein, starting in 1834. The takeoff stage of economic development came with the railroad revolution in the 1840s, which opened up new markets for local products, created a pool of middle managers, increased the demand for engineers, architects and skilled machinists, and stimulated investments in coal and iron. The political decisions about the economy of Prussia (and after 1871 all Germany) were largely controlled by a coalition of "rye and iron", that is the Junker landowners of the east and the heavy industry of the west. The north German states were for the most part richer in natural resources than the southern states. They had vast agricultural tracts from Schleswig-Holstein in the west through Prussia in the east. They also had coal and iron in the Ruhr Valley. Through the practice of primogeniture, widely followed in northern Germany, large estates and fortunes grew. So did close relations between the owners and local as well as national governments. The south German states were relatively poor in natural resources and those Germans therefore engaged more often in small economic enterprises. They also had no primogeniture rule but subdivided the land among several offspring, leading those offspring to remain in their native towns but not fully able to support themselves from their small parcels of land. The south German states, therefore, fostered cottage industries, crafts, and a more independent and self-reliant spirit less closely linked to the government. The first important mines appeared in the 1750s, in the valleys of the rivers Ruhr, Inde and Wurm where coal seams outcropped and horizontal adit mining was possible. In 1782 the Krupp family began operations near Essen. After 1815 entrepreneurs in the Ruhr Area, which then became part of Prussia, took advantage of the tariff zone (Zollverein) to open new mines and associated iron smelters. New railroads were built by British engineers around 1850. Numerous small industrial centres sprang up, focused on ironworks, using local coal. The iron and steel works typically bought mines and erected coking ovens to supply their own requirements in coke and gas. These integrated coal-iron firms ("Huettenzechen") became numerous after 1854; after 1900 they became mixed firms called "Konzern." The output of an average mine in 1850 was about 8,500 short tons; its employment about 64. By 1900, this output had risen to 280,000 and employment to about 1,400. Total Ruhr coal output rose from 2.0 million short tons in 1850 to 22 in 1880, 60 in 1900, and 114 in 1913, on the verge of war. In 1932 output was down to 73 million short tons, growing to 130 in 1940. Output peaked in 1957 (at 123 million), declining to 78 million short tons in 1974. End of 2010 five coal mines were producing in Germany. The miners in the Ruhr region were divided by ethnicity (Germans and Poles) and religion (Protestants and Catholics). Mobility in and out of the mining camps to nearby industrial areas was high. The miners split into several unions, with an affiliation to a political party. As a result, the socialist union (affiliated with the Social Democratic Party) competed with Catholic and Communist unions until 1933, when the Nazis took over all of them. After 1945 the socialists came to the fore. Banks and cartels German banks played central roles in financing German industry. Different banks formed cartels in different industries. Cartel contracts were accepted as legal and binding by German courts although they were held to be illegal in Britain and the United States. The process of cartelization began slowly, but the cartel movement took hold after 1873 in the economic depression that followed the postunification speculative bubble. It began in heavy industry and spread throughout other industries. By 1900 there were 275 cartels in operation; by 1908, over 500. By some estimates, different cartel arrangements may have numbered in the thousands at different times, but many German companies stayed outside the cartels because they did not welcome the restrictions that membership imposed. The government played a powerful role in the industrialization of the German Empire founded by Otto von Bismarck in 1871 during a period known as the Second Industrial Revolution. It supported not only heavy industry but also crafts and trades because it wanted to maintain prosperity in all parts of the empire. Even where the national government did not act, the highly autonomous regional and local governments supported their own industries. Each state tried to be as self-sufficient as possible. Despite the several ups and downs of prosperity and depression that marked the first decades of the German Empire, the ultimate wealth of the empire proved immense. German aristocrats, landowners, bankers, and producers created what might be termed the first German economic miracle, the turn-of-the-century surge in German industry and commerce during which bankers, industrialists, mercantilists, the military, and the monarchy joined forces. Class and the welfare state Germany's middle class, based in the cities, grew exponentially, but it never gained the political power it had in France, Britain or the United States. The Association of German Women's Organizations (BDF) was established in 1894 to encompass the proliferating women's organizations that had sprung up since the 1860s. From the beginning the BDF was a bourgeois organization, its members working toward equality with men in such areas as education, financial opportunities, and political life. Working-class women were not welcome; they were organized by the Socialists. Bismarck built on a tradition of welfare programs in Prussia and Saxony that began as early as in the 1840s. In the 1880s he introduced old age pensions, accident insurance, medical care and unemployment insurance that formed the basis of the modern European welfare state. His paternalistic programs won the support of German industry because its goals were to win the support of the working classes for the Empire and reduce the outflow of immigrants to America, where wages were higher, but welfare did not exist. Bismarck further won the support of both industry and skilled workers by his high tariff policies, which protected profits and wages from American competition, although they alienated the liberal intellectuals who wanted free trade. Political disunity of three dozen states and a pervasive conservatism made it difficult to build railways in the 1830s. However, by the 1840s, trunk lines did link the major cities; each German state was responsible for the lines within its own borders. Economist Friedrich List summed up the advantages to be derived from the development of the railway system in 1841: - as a means of national defence, it facilitates the concentration, distribution and direction of the army. - It is a means to the improvement of the culture of the nation…. It brings talent, knowledge and skill of every kind readily to market. - It secures the community against dearth and famine, and against excessive fluctuation in the prices of the necessaries of life. - It promotes the spirit of the nation, as it has a tendency to destroy the Philistine spirit arising from isolation and provincial prejudice and vanity. It binds nations by ligaments, and promotes an interchange of food and of commodities, thus making it feel to be a unit. The iron rails become a nerve system, which, on the one hand, strengthens public opinion, and, on the other hand, strengthens the power of the state for police and governmental purposes. Lacking a technological base at first, the Germans imported their engineering and hardware from Britain, but quickly learned the skills needed to operate and expand the railways. In many cities, the new railway shops were the centres of technological awareness and training, so that by 1850, Germany was self-sufficient in meeting the demands of railroad construction, and the railways were a major impetus for the growth of the new steel industry. Observers found that even as late as 1890, their engineering was inferior to Britain’s. However, German unification in 1870 stimulated consolidation, nationalisation into state-owned companies, and further rapid growth. Unlike the situation in France, the goal was support of industrialisation, and so heavy lines crisscrossed the Ruhr and other industrial districts, and provided good connections to the major ports of Hamburg and Bremen. By 1880, Germany had 9,400 locomotives pulling 43,000 passengers and 30,000 tons of freight, and pulled ahead of France. Perkins (1981) argues that more important than Bismarck's new tariff on imported grain was the introduction of the sugar beet as a primary crop. Farmers quickly abandoned traditional, inefficient practices for modern new methods, including use of new fertilizers and new tools. The knowledge and tools gained from the intensive farming of sugar and other root crops made Germany the most efficient agricultural producer in Europe by 1914. Even so, farms were small in size, and women did much of the field work. An unintended consequence was the increased dependence on migratory, especially foreign, labor. The economy continued to industrialize and urbanize, with heavy industry (coal and steel especially) becoming important in the Ruhr, and manufacturing growing in the cities, the Ruhr, and Silesia. Based on its leadership in chemical research in the universities and industrial laboratories, Germany became dominant in the world's chemical industry in the late 19th century. Big businesses such as BASF and Bayer led the way in their production and distribution of artificial dyes and pharmaceuticals during the Wilhelmine era, leading to the German monopolisation of the global chemicals market at 90 percent of the entire share of international volumes of trade in chemical products by 1914. Germany became Europe's leading steel-producing nations in the late 19th century, thanks in large part to the protection from American and British competition afforded by tariffs and cartels. The leading firm was "Friedrich Krupp AG Hoesch-Krupp" run by the Krupp family The "German Steel Federation" was established in 1874. The merger of four major firms into the Vereinigte Stahlwerke (United Steel Works) in 1926 was modeled on the U.S. Steel corporation in the U.S. The goal was to move beyond the limitations of the old cartel system by incorporating advances simultaneously inside a single corporation. The new company emphasized rationalization of management structures and modernization of the technology; it employed a multi-divisional structure and used return on investment as its measure of success. By 1913 American and German exports dominated the world steel market, as Britain slipped to third place. In machinery, iron and steel and other industries, German firms avoided cut-throat competition and instead relied on trade associations. Germany was a world leader because of its prevailing "corporatist mentality", its strong bureaucratic tradition, and the encouragement of the government. These associations regulated competition and allowed small firms to function in the shadow of much larger companies. First World War British economist John Maynard Keynes denounced the 1919 Treaty of Versailles as ruinous to German and global prosperity. In his book The Economic Consequences of the Peace. Keynes said the Treaty was a "Carthaginian peace", a misguided attempt to destroy Germany on behalf of French revanchism, rather than to follow the fairer principles for a lasting peace set out in President Woodrow Wilson's Fourteen Points, which Germany had accepted at the armistice. Keynes argued the sums being asked of Germany in reparations were many times more than it was possible for Germany to pay, and that these would produce drastic instability. French economist Étienne Mantoux disputed that analysis in The Carthaginian Peace, or the Economic Consequences of Mr. Keynes (1946). More recently economists have argued that the restriction of Germany to a small army in the 1920s saved it so much money it could afford the reparations payments. In reality, the total German Reparation payments actually made were far smaller than anyone expected. The total came to 20 billion German gold marks, worth about $5 billion US dollars or £1 billion British pounds. German reparations payments ended in 1931. The war and the treaty were followed by the Hyper-inflation of the early 1920s that wreaked havoc on Germany's social structure and political stability. During that inflation, the value of the nation's currency, the Papiermark, collapsed from 8.9 per US$1 in 1918 to 4.2 trillion per US$1 by November 1923. Prosperity reigned 1923–29, supported by large bank loans from New York. The Great Depression struck Germany hard, starting in late 1929. There were no new American loans. Unemployment soared, especially in larger cities, fueling extremism and violence on the far right and far left, as the centre of the political spectrum weakened. Germany had paid about one-eighth of its war reparations when they were suspended in 1932 by the Lausanne Conference of 1932. The failure of major banks in Germany and Austria in 1931 worsened the worldwide banking crisis. During the Hitler era (1933–45), the economy developed a hothouse prosperity, supported with high government subsidies to those sectors that tended to give Germany military power and economic autarky, that is, economic independence from the global economy. During the war itself the German economy was sustained by the exploitation of conquered territories and people. With the loss of the war, the country entered into the period known as Stunde Null ("Zero Hour"), when Germany lay in ruins and the society had to be rebuilt from scratch. Post-World War II The first several years after World War II were years of bitter penury for the Germans. Seven million forced laborers left for their own land, but about 14 million Germans came in from the East, living for years in dismal camps. It took nearly a decade for all the German POWs to return. In the West, farm production fell, food supplies were cut off from eastern Germany (controlled by the Soviets) and food shipments extorted from conquered lands ended. The standard of living fell to levels not seen in a century, and food was always in short supply. High inflation made savings (and debts) lose 99% of their value, while the black market distorted the economy. In the East, the Soviets crushed dissent and imposed another police state, often employing ex-Nazis in the dreaded Stasi. The Soviets extracted about 23% of the East German GNP for reparations, while in the West reparations were a minor factor. The man who took full advantage of Germany's postwar opportunity was Ludwig Erhard, who was determined to shape a new and different kind of German economy. He was given his chance by United States officials, who found him working in Nuremberg and who saw that many of his ideas coincided with their own. Erhard abolished the Reichsmark and the created a new currency, the Deutsche Mark, on 21 June 1948, with the concurrence of the Western Allies but also taking advantage of the opportunity to abolish most Nazi and occupation rules and regulations. It established the foundations of the West German economy and of the West German state. After 1950, Germany overtook Britain in comparative productivity levels for the whole economy, primarily as a result of trends in services rather than trends in industry. The Marshall Plan was eagerly adopted in Germany as a way to modernize business procedures and utilize the best practices, while these changes were resisted in Britain. Britain's historic lead in productivity of its services sector was based on external economies of scale in a highly urbanized economy with an international orientation. On the other hand, the low productivity in Germany was caused by the underdevelopment of services generally, especially in rural areas that comprised a much larger sector. As German farm employment declined sharply after 1950 thanks to mechanization, catching-up occurred in services. This process was aided by a sharp increase in human and physical capital accumulation, a pro-growth government policy, and the effective utilization of the education sector to create a more productive work force. Social market economy The Germans proudly label their economy a "soziale Marktwirtschaft," or "social market economy," to show that the system as it has developed after World War II has both a material and a social—or human—dimension. They stress the importance of the term "market" because after the Nazi experience they wanted an economy free of state intervention and domination. The only state role in the new West German economy was to protect the competitive environment from monopolistic or oligopolistic tendencies—including its own. The term "social" is stressed because West Germans wanted an economy that would not only help the wealthy but also care for the workers and others who might not prove able to cope with the strenuous competitive demands of a market economy. The term "social" was chosen rather than "socialist" to distinguish their system from those in which the state claimed the right to direct the economy or to intervene in it. Beyond these principles of the social market economy, but linked to it, comes a more traditional German concept, that of Ordnung, which can be directly translated to mean order but which really means an economy, society, and policy that are structured but not dictatorial. The founders of the social market economy insisted that Denken in Ordnungen—to think in terms of systems of order—was essential. They also spoke of Ordoliberalism because the essence of the concept is that this must be a freely chosen order, not a command order. Over time, the term "social" in the social market economy began to take on a life of its own. It moved the West German economy toward an extensive social welfare system that has become one of the most expensive in the world. Moreover, the West German federal government and the states (Länder ; sing., Land ) began to compensate for irregularities in economic cycles and for shifts in world production by beginning to shelter and support some sectors and industries. In an even greater departure from the Erhard tradition, the government became an instrument for the preservation of existing industries rather than a force for renewal. In the 1970s, the state assumed an ever more important role in the economy. During the 1980s, Chancellor Helmut Kohl tried to reduce that state role, and he succeeded in part, but German unification again compelled the German government to assume a stronger role in the economy. Thus, the contradiction between the terms "social" and "market" has remained an element for debate in Germany. Given the internal contradiction in its philosophy, the German economy is both conservative and dynamic. It is conservative in the sense that it draws on the part of the German tradition that envisages some state role in the economy and a cautious attitude toward investment and risk-taking. It is dynamic in the sense that it is directed toward growth—even if that growth may be slow and steady rather than spectacular. It tries to combine the virtues of a market system with the virtues of a social welfare system. Economic miracle and beyond The economic reforms and the new West German system received powerful support from a number of sources: investment funds under the European Recovery Program, more commonly known as the Marshall Plan; the stimulus to German industry provided by the diversion of other Western resources for Korean War production; and the German readiness to work hard for low wages until productivity had risen. But the essential component of success was the revival of confidence brought on by Erhard's reforms and by the new currency. The West German boom that began in 1950 was truly memorable. The growth rate of industrial production was 25.0 percent in 1950 and 18.1 percent in 1951. Growth continued at a high rate for most of the 1950s, despite occasional slowdowns. By 1960 industrial production had risen to two-and-one-half times the level of 1950 and far beyond any that the Nazis had reached during the 1930s in all of Germany. GDP rose by two-thirds during the same decade. The number of persons employed rose from 13.8 million in 1950 to 19.8 million in 1960, and the unemployment rate fell from 10.3 percent to 1.2 percent. Labor also benefited in due course from the boom. Although wage demands and pay increases had been modest at first, wages and salaries rose over 80 percent between 1949 and 1955, catching up with growth. West German social programs were given a considerable boost in 1957, just before a national election, when the government decided to initiate a number of social programs and to expand others. In 1957 West Germany gained a new central bank, the Deutsche Bundesbank, generally called simply the Bundesbank, which succeeded the Bank deutscher Länder and was given much more authority over monetary policy. That year also saw the establishment of the Bundeskartellamt (Federal Cartel Office), designed to prevent the return of German monopolies and cartels. Six years later, in 1963, the Bundestag, the lower house of Germany's parliament, at Erhard's urging established the Council of Economic Experts to provide objective evaluations on which to base German economic policy. The West German economy did not grow as fast or as consistently in the 1960s as it had during the 1950s, in part because such a torrid pace could not be sustained, in part because the supply of fresh labor from East Germany was cut off by the Berlin Wall, built in 1961, and in part because the Bundesbank became disturbed about potential overheating and moved several times to slow the pace of growth. Erhard, who had succeeded Konrad Adenauer as chancellor, was voted out of office in December 1966, largely—although not entirely—because of the economic problems of the Federal Republic. He was replaced by the Grand Coalition consisting of the Christian Democratic Union (Christlich Demokratische Union—CDU), its sister party the Christian Social Union (Christlich-Soziale Union—CSU), and the Social Democratic Party of Germany (Sozialdemokratische Partei Deutschlands—SPD) under Chancellor Kurt Georg Kiesinger of the CDU. Under the pressure of the slowdown, the new West German Grand Coalition government abandoned Erhard's broad laissez-faire orientation. The new minister for economics, Karl Schiller, argued strongly for legislation that would give the federal government and his ministry greater authority to guide economic policy. In 1967 the Bundestag passed the Law for Promoting Stability and Growth, known as the Magna Carta of medium-term economic management. That law, which remains in effect although never again applied as energetically as in Schiller's time, provided for coordination of federal, Land, and local budget plans in order to give fiscal policy a stronger impact. The law also set a number of optimistic targets for the four basic standards by which West German economic success was henceforth to be measured: currency stability, economic growth, employment levels, and trade balance. Those standards became popularly known as the magisches Viereck, the "magic rectangle" or the "magic polygon." Schiller followed a different concept from Erhard's. He was one of the rare German Keynesians, and he brought to his new tasks the unshakable conviction that government had both the obligation and the capacity to shape economic trends and to smooth out and even eliminate the business cycle. Schiller's chosen formula was Globalsteuerung, or global guidance, a process by which government would not intervene in the details of the economy but would establish broad guidelines that would foster uninterrupted noninflationary growth. Schiller's success in the Grand Coalition helped to give the SPD an electoral victory in 1969 and a chance to form a new coalition government with the Free Democratic Party (Freie Demokratische Partei—FDP) under Willy Brandt. The SPD-FDP coalition expanded the West German social security system, substantially increasing the size and cost of the social budget. Social program costs grew by over 10 percent a year during much of the 1970s, introducing into the budget an unalterable obligation that reduced fiscal flexibility (although Schiller and other Keynesians believed that it would have an anticyclical effect). This came back to haunt Schiller as well as every German government since then. Schiller himself had to resign in 1972 when the West German and global economies were in a downturn and when all his ideas did not seem able to revive West German prosperity. Willy Brandt himself resigned two years later. Helmut Schmidt, Brandt's successor, was intensely interested in economics but also faced great problems, including the dramatic upsurge in oil prices of 1973-74. West Germany's GDP in 1975 fell by 1.4 percent (in constant prices), the first time since the founding of the FRG that it had fallen so sharply. The West German trade balance also fell as global demand declined and as the terms of trade deteriorated because of the rise in petroleum prices. By 1976 the worst was over. West German growth resumed, and the inflation rate began to decline. Although neither reached the favorable levels that had come to be taken for granted during the 1950s and early 1960s, they were accepted as tolerable after the turbulence of the previous years. Schmidt began to be known as a Macher (achiever), and the government won reelection in 1976. Schmidt's success led him and his party to claim that they had built Modell Deutschland (the German model). But the economy again turned down and, despite efforts to stimulate growth by government deficits, failed to revive quickly. It was only by mid-1978 that Schmidt and the Bundesbank were able to bring the economy into balance. After that, the economy continued expanding through 1979 and much of 1980, helping Schmidt win reelection in 1980. But the upturn proved to be uneven and unrewarding, as the problems of the mid-1970s rapidly returned. By early 1981, Schmidt faced the worst possible situation: growth fell and unemployment rose, but inflation did not abate. By late 1982, Schmidt's coalition government collapsed as the FDP withdrew to join a coalition led by Helmut Kohl, the leader of the CDU/CSU. He began to direct what was termed the Wende (West Germany) (turning or reversal). The government proceeded to implement new policies to reduce the government role in the economy and within a year won a popular vote in support of the new course. Within its broad policy, the new government had several main objectives: to reduce the federal deficit by cutting expenditures as well as taxes, to reduce government restrictions and regulations, and to improve the flexibility and performance of the labor market. The government also carried through a series of privatization measures, selling almost DM10 billion (for value of the deutsche mark—see Glossary) in shares of such diverse state-owned institutions as VEBA, VIAG, Volkswagen, Lufthansa, and Salzgitter. Through all these steps, the state role in the West German economy declined from 52 percent to 46 percent of GDP between 1982 and 1990, according to Bundesbank statistics. Although the policies of the Wende changed the mood of the West German economy and reinstalled a measure of confidence, progress came unevenly and haltingly. During most of the 1980s, the figures on growth and inflation improved but slowly, and the figures on unemployment barely moved at all. There was little job growth until the end of the decade. When the statistics did change, however, even modestly, it was at least in the right direction. Nonetheless, it also remained true that West German growth did not again reach the levels that it had attained in the early years of the Federal Republic. There had been a decline in the growth rate since the 1950s, an upturn in unemployment since the 1960s, and a gradual increase in inflation except during or after a severe downturn. Global economic statistics also showed a decline in West German output and vitality. They showed that the West German share of total world production had grown from 6.6 percent in 1965 to 7.9 percent by 1975. Twelve years later, in 1987, however, it had fallen to 7.4 percent, largely because of the more rapid growth of Japan and other Asian states. Even adding the estimated GDP of the former East Germany at its peak before unification would not have brought the all-German share above 8.2 percent by 1989 and would leave all of Germany with barely a greater share of world production than West Germany alone had reached fifteen years earlier. It was only in the late 1980s that West Germany's economy finally began to grow more rapidly. The growth rate for West German GDP rose to 3.7 percent in 1988 and 3.6 percent in 1989, the highest levels of the decade. The unemployment rate also fell to 7.6 percent in 1989, despite an influx of workers from abroad. Thus, the results of the late 1980s appeared to vindicate the West German supply-side revolution. Tax rate reductions had led to greater vitality and revenues. Although the cumulative public-sector deficit had gone above the DM1 trillion level, the public sector was growing more slowly than before. The year 1989 was the last year of the West German economy as a separate and separable institution. From 1990 the positive and negative distortions generated by German reunification set in, and the West German economy began to reorient itself toward economic and political union with what had been East Germany. The economy turned gradually and massively from its primarily West European and global orientation toward an increasingly intense concentration on the requirements and the opportunities of unification. German reunification and its aftermath Germany invested over $2 trillion marks in the rehabilitation of the former East Germany helping it to transition to a market economy, and cleaning up the environmental degradation. By 2011 the results were mixed, with slow economic development in the East, in sharp contrast to the rapid economic growth in both west and southern Germany. Unemployment was much higher in the East, often over 15%. Economists Snower and Merkl (2006) suggests that the malaise was prolonged by all the social and economic help from the German government, pointing especially to bargaining by proxy, high unemployment benefits and welfare entitlements, and generous job security provisions. The old industrial centers of the Rhineland and North Germany lagged as well, as the coal and steel industries faded in importance. The economic policies were heavily oriented toward the world market, and the export sector continued very strong. - "Germany". CIA World Factbook. Central Intelligence Agency. 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Routledge. pp. 32–34. - Richard Tilly, "Germany: 1815-1870" in Rondo Cameron, ed. Banking in the Early Stages of Industrialization: A Study in Comparative Economic History (Oxford University Press, 1967), pages 151-182 - Cornelius Torp, "The "Coalition of 'Rye and Iron'" under the Pressure of Globalization: A Reinterpretation of Germany's Political Economy before 1914," Central European History Sept 2010, Vol. 43 Issue 3, pp 401-427 - Griffin, Emma. "Why was Britain first? The Industrial revolution in global context". Short History of the British Industrial Revolution. Retrieved 6 February 2013. - Pounds (1952) - Stefan Llafur Berger, "Working-Class Culture and the Labour Movement in the South Wales and the Ruhr Coalfields, 1850-2000: A Comparison," Journal of Welsh Labour History/Cylchgrawn Hanes Llafur Cymru (2001) 8#2 pp 5-40. - Eda Sagarra, A Social History of Germany 1648-1914 (2002) - E. P. Hennock, The Origin of the Welfare State in England and Germany, 1850–1914: Social Policies Compared (2007); Hermann Beck, Origins of the Authoritarian Welfare State in Prussia, 1815-1870 (1995) - Elaine Glovka Spencer, "Rules of the Ruhr: Leadership and Authority in German Big Business Before 1914," Business History Review, Spring 1979, Vol. 53 Issue 1, pp 40-64; Ivo N. Lambi, "The Protectionist Interests of the German Iron and Steel Industry, 1873-1879," Journal of Economic History, March 1962, Vol. 22 Issue 1, pp 59-70 - List quoted in John J. Lalor, ed. Cyclopædia of Political Science (1881) 3:118 online; see Thomas Nipperdey, Germany from Napoleon to Bismarck (1996) p 165 - Allan Mitchell, Great Train Race: Railways and the Franco-German Rivalry, 1815-1914 (2000) - J.A. Perkins, "The Agricultural Revolution in Germany 1850–1914," Journal of European Economic History, Spring 1981, Vol. 10 Issue 1, pp 71-119 - Cornelius Torp, "The Great Transformation: German Economy and Society 1850-1914", in Helmut Walser Smith (ed.), The Oxford Handbook of Modern German History (2011), pp. 347-8 - Steven B. Webb, "Tariffs, Cartels, Technology, and Growth in the German Steel Industry, 1879 to 1914," Journal of Economic History Vol. 40, No. 2 (Jun., 1980), pp. 309-330 in JSTOR - Harold James, Krupp: A History of the Legendary German Firm (Princeton U.P. 2012) - "The German Steel Federation". WV Stahl. Archived from the original on 2007-01-29. Retrieved 2007-04-26. - Alfred Reckendrees, "From Cartel Regulation to Monopolistic Control? The Founding of the German 'Steel Trust' in 1926 and its Effect on Market Regulation," Business History, (July 2003) 45#3 pp 22-51, - Robert C. Allen, "International Competition in Iron and Steel, 1850-1913, Journal of Economic History, (Dec 1979) 39#4 pp 911-37 in JSTOR - Gerald D. Feldman and Ulrich Nocken, "Trade Associations and Economic Power: Interest Group Development in the German Iron and Steel and Machine Building Industries, 1900-1933" Business History Review, (Winter 1975), 49#4 pp 413-45 in JSTOR - Feldman, Gerald D. "The Political and Social Foundations of Germany's Economic Mobilization, 1914-1916," Armed Forces & Society (1976) 3#1 pp 121-145. online - Antony. Lentin, "Germany: a New Carthage?," History Today (Jan. 2012) 62#1 pp 20-27 - Keynes (1919). The Economic Consequences of the Peace. Ch VI. - Max Hantke and Mark Spoerer, "The imposed gift of Versailles: the fiscal effects of restricting the size of Germany's armed forces, 1924-9," Economic History Review (2010) 63#4 pp 849-864. online - Sally. 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Optimizing the German Workforce: Labor Administration From Bismarck to the Economic Miracle (Berghahn Books; 2010) 276 pages; studies continuities in German governments' efforts to create a skilled labor force across the disparate imperial, Weimar, Nazi, and postwar regimes. - Milward, Alan S. and S. B. Saul. The Development of the Economies of Continental Europe: 1850-1914 (1977) pp 17–70 - Milward, Alan S. and S. B. Saul. The Economic Development of Continental Europe 1780-1870 (1973) - Overy, R. J. The Nazi Economic Recovery 1932-1938 (1996) excerpt and text search - Overy, R. J. War and Economy in the Third Reich (1994) - Perkins, J. A. "Dualism in German Agrarian Historiography, Comparative Studies in Society & History, Apr 1986, Vol. 28 Issue 2, pp 287–330, compares large landholdings in the territories east of the Elbe river, and the West-Elbian small-scale agriculture. - Pierenkemper, T., and R. Tilly, The German Economy during the Nineteenth Century (2004) - Sagarra, Eda. 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In the middle of the eighteenth century, Jews living in German territories were just beginning to feel the effects of the political, social and intellectual changes that would soon be recognized as the hallmarks of the modern world. Until this period, Jewish communities had been constituted as distinct and autonomous social, religious and legal entities within an essentially feudal social organization. The Jews were subjected to the will of the rulers of individual German states, who imposed onerous regulations, taxes and restrictions on their ability to marry and settle where they chose. Distinguished from the rest of the population by religious traditions and family structure, Jews lived under the authority of the Jewish community, wholly separate from the non-Jewish population. In the early 1780s, however, Enlightenment thinkers began to call for an end to the discrimination against Jews in Prussia and Austria. Most important among these voices was the high-ranking Prussian government official Christian Wilhelm Dohm (1751–1820), who argued that Jews be granted the same civil rights as those accorded to non-Jewish citizens. In his essay “On the Civil Betterment of the Jews” (1781) Dohm explained Jews’ moral “depravity” as the result of centuries of oppression. Only with the elimination of the oppressive conditions that produced their allegedly defective character, Dohm argued, would Jews be able to gradually overcome their “disabilities” and prove themselves to be useful citizens. These changes in the intellectual realm coincided with broader social and political realignments of the period. In its attempt to centralize its power and eliminate intermediate corporate bodies such as guilds and estates, the emerging absolutist state also sought the dissolution of the autonomous Jewish community and the integration of its members into the larger social body. Thus the development of more modern social and political structures for German Jews arose as much from larger external factors as from a desire to address the particular situation of the Jews. Indeed, the practical reforms that were introduced by Emperor Joseph II in his “Tolerance Decree” of 1781 probably exerted a more immediate impact on the situation of Jews in parts of the German-speaking territories than did Enlightenment thought itself. Joseph II’s decree enacted the first legal measures to reduce legal restrictions on the Jewish population in parts of the Habsburg Empire. Despite the beginnings of a more consolidated state authority, Jewish legal status still varied within the territories of the German Empire until the unification of Germany in 1871. Simultaneous with the political developments that gradually began to erode the legal barriers separating Jew from non-Jew were important changes that also took place within Jewish society itself. Influenced by the spirit of the Enlightenment, Jewish intellectuals began a new critical engagement with Jewish tradition and, in so doing, created a new cultural, social and intellectual framework and helped bring forth an invigorated public sphere that challenged the authority of the official Jewish community. A new intellectual elite emerged, distinct from the rabbinate and its influence. The combined impact of the centralizing absolutist state and the emergence of the European and Jewish Enlightenments marked the beginning of a change in the legal status of the Jews that would extend over more than a hundred-year period. Yet the progress of Jewish Emancipation in different territories was anything but linear. During periods of political liberalization, progress toward equal rights proceeded apace, but the process was set back during periods of conservative counterreaction. Because of the protracted nature of the struggle for Jewish Emancipation, and the twin efforts to win both legal and social acceptance, the uneven development of Emancipation was a central defining experience for German Jews. Yet historians have generally treated the Emancipation of the Jews as an event of universal significance for German Jewry without paying significant attention to the gendered aspects of its unfolding. After a long and often painful process, Jewish men did finally achieve full political and civil rights with the unification of Germany in 1871. In principle, if not entirely in practice, this removed most remaining legal disabilities that had prevented the full integration of male Jews into German society. But at the time of Emancipation, Jewish women—like women in general—received no such rights and remained unable to vote until 1918. In fact, as men looked toward an era of increasing liberalization, women were politically disenfranchised as the result of a law that was in effect until 1908, banning women from joining political organizations. Although German women were citizens, their status was ultimately determined by the citizenship of their father or husband. The status of East European Jewish women was even more precarious, since many immigrant women were not permitted to become citizens at all. Within the Jewish community, Jewish women had to wait even longer to gain a voice and a vote. In Germany, Jewish women suffered the double indignity of sexism and antisemitism, with second-class status imposed both inside and outside the Jewish community. If Emancipation affected Jewish men and women in distinct ways, the pace and extent to which Jews adapted themselves to the demands of German society also differed according to gender-specific patterns. Because social acceptance was made contingent upon the acquisition of the basic customs, behaviors and values of German society, Jewish men and women took different paths toward, and found different means for, becoming at once fully German and distinctly Jewish. Jewish men tended to adapt to the demands of middle-class society by abandoning public religious behaviors, including the observance of Jewish dietary laws and the prohibition of work on the Sabbath. They also concerned themselves less with Jewish learning and worship than with secular education, which they pursued with unparalleled enthusiasm. For women, the road to acculturation led to the formation of new roles inside and outside the home. Changes in family structure and employment patterns led to adjustments in the gender division of labor between the domestic and public spheres, and these changes in the family, religion and labor, in turn, affected the construction of gender roles within the Jewish community, and gender relations as a whole. Even the category “Jewish woman” was infused with new meanings that accorded with middle-class norms and ideals for the bourgeois German woman. In fulfilling their newly defined “woman’s nature,” Jewish women created a proliferation of voluntary associations, involved themselves in non-Jewish associations, and pioneered the field of social work. The path to becoming “German” and “German Jewish” thus proved to be profoundly gendered. One of the earliest examples of a specifically female experience of the Enlightenment can be found in the Berlin salons of the late eighteenth century. Although the German Jewish Enlightenment is usually associated with the intellectual circle around Moses Mendelssohn and its literary output, most historical literature has treated the Enlightenment as an intellectual and socio-cultural phenomenon that has almost exclusively involved men. Yet during the last two decades of the eighteenth century, even as Mendelssohn was evolving his philosophical reformulation of Judaism, a small group of young women from Berlin’s small but influential Jewish upper class crafted a place for themselves at the very center of the city’s social and intellectual life. Jewish salionières, most notably Rachel Levin Varnhagen and Dorothea Schlegel Mendelssohn, hosted social and intellectual gatherings in their homes that brought together Jews and non-Jews, noblemen and commoners to socialize and exchange ideas. Creating a cultural space unprecedented in its openness to Jews and women, these salons appear to have existed only for a brief historical moment as if outside the normal social constraints that enforced the hierarchical organization of society around the axes of gender, class and religion. Because a substantial number of these women converted to Christianity and entered into (often second) marriages with non-Jews, Jewish historians have sometimes been quick to condemn them for the betrayal of their people and faith. So traitorous were they, concluded Heinrich Graetz, that “these talented but sinful Jewish women did Judaism a service by becoming Christians” (Lowenstein, 109). Indeed, for their contemporaries, as well as for many historians, these women represented the embodiment of a larger set of social problems afflicting Berlin Jewish society. Whatever the exact mix of their motives for conversion and intermarriage—an ascent in social status, the promise of companionate marriage or liberation from their patriarchal families—the fact that contemporary observers and later historians held these boundary-crossing women responsible for the most visible symptoms of modern social change suggests the extent to which the transition from traditional Jewish society to modern Judaism was represented through the language of gender. Among the earliest and most important trajectories for the progress of German Jewish acculturation was the modernization of Judaism. Although religious modernization by no means did away with gender hierarchy, it nevertheless altered gender expectations and gender roles, as well as popular forms of religious practice. Within traditional Judaism, those aspects of religious practice that had historically been invested with the greatest value were organized hierarchically along clearly gendered lines: the study of Torah and public prayer formed the religious centerpiece of Jewish life, and these acts were accessible only to men. Women’s religious activity tended to be less structured and more personal and took place primarily in the family, focusing on religious aspects of home life, the observance of Sabbath and holidays, and the maintenance of dietary and family purity laws. Though accorded importance within Judaism, these “domestic” practices of Judaism were lower in prestige than the more public practice of Judaism dominated by men. The liberal religious reform movement, which has garnered so much attention in the historical literature, did not, however, fundamentally transform the role of women in Judaism. Reformers seeking to modernize Judaism in accordance with Enlightenment ideals and middle-class behavioral and aesthetic visions endeavored to make prayer services more attractive to women as well as men by changing the language of prayer from Hebrew to German and replacing Hebrew excurses on the law with uplifting preaching in German that was modeled on Protestant worship services. Equally important, greater attention was paid to women’s religious education, primarily through the inclusion of women in the newly introduced ritual of confirmation. There was even some discussion at the 1846 Rabbinic Assembly in Breslau of far-reaching changes that would have granted women greater religious equality. Yet when it came to practice, the nineteenth-century Reform movement failed to eliminate many of the traditional religious restrictions that kept women in a subordinate status. In the synagogue, women could still neither be counted in a prayer quorum nor called to the Torah, and they often remained seated apart and comfortably out of view in the women’s gallery. Despite pronouncements against the segregation of women, religious reformers ultimately made few substantial improvements in women’s status. In addition to religious reform, religious modernization also includes those religious and cultural changes that resulted from the increased participation of women in religious associations outside the home and the formal sphere of the Jewish community. Indeed, it may well have been this phenomenon, more than religious reform itself, which affected gender relations more broadly and contributed to a more substantive reconfiguration of the traditional Jewish gender order. Beginning in the late eighteenth century, middle-class women began to create charitable associations, such as sick care and self-help societies, that mirrored both the form and content of male associations. According to the historian Maria Benjamin Baader, these new female voluntary organizations were made possible with the declining emphasis on traditional male learning that had once marked women as marginal. Expressing both Jewish and bourgeois values, women’s activity in this realm would, by the early twentieth century, lead into new professional opportunities as well as to the production of new forms of Jewish religious and ethnic expression. Women’s activities in voluntary organizations, in turn, were linked to broader cultural shifts in the German bourgeoisie. Thus, in addition to being inflected by gender, it is important to note that the process of becoming a German Jewish woman was also affected by class status. Whereas on the eve of the modern era the majority of German Jews were poor, as they entered German society over the course of the nineteenth century, Jews aspired to join the class that was most suited to their skills in trade and commerce: the middle-class. Jews quickly embraced the ideal of the educated middle class that made culture, rather than birth, the defining character of class. While middle-class status for men was to be achieved through self-improvement and education (Bildung), the most important determinant of middle class respectability for women and for their husbands was her status as full-time “priestess of the home.” Paradoxically, an important measure of specifically Jewish and middle class acculturation was a new form of family-centered Judaism that arose out of the strong emphasis placed on the family in bourgeois culture on the one hand, and the decline in traditional Jewish religious practice on the other. The nineteenth-century bourgeois ideal for the family was a prescriptive model based on a rigid gender-based division of labor that delimited women’s activities to the domestic sphere and men’s activity to the “public” arena. As an ideal, it quickly eclipsed the typical structure of premodern Jewish families where the boundaries between public and private remained more fluid. Of course, not all Jewish families could afford to imitate this model, since lower middle-class and working class families often had to rely on the work and wages of children and wives for the family economy. But the power of this construct as a universal model for Jewish family life is perhaps most evidenced by the fact that, since the mid-nineteenth century, the bourgeois family type has been viewed as the “traditional Jewish family.” By the time the German states were joined in a federal system within the new German Empire in 1871, most of Germany’s Jews could proudly display their middle class status by pointing to their family life. Indeed, the research of Marion Kaplan has demonstrated how Jewish women managed the double task of transmitting the values and behaviors of the German bourgeoisie while helping to shape the Jewish identity of their children. Jewish women made sure their children learned the German classics and, at the same time, organized the observance of holidays, family gatherings and the religious and moral education of the children. Illustrating the family’s crucial role in the acculturation of German Jews, Kaplan’s research also suggests the extent to which the home was gradually being recast as the primary site for the transmission of Judaism. With the declining appeal of formal religious practice and institutions, including the synagogue, the Jewish mother, according to historian Jacob Toury, was expected to become the “protector of a new system of Jewish domestic culture” (Maurer, 147). Although some historians suggest that Jewish men abandoned religious ritual and practice more quickly than women, by mid-century Jewish community leaders nevertheless began holding women increasingly accountable for assimilation, conversion and intermarriage—in short, for the decline of Judaism. This was the case despite the fact that the intermarriage and conversion rates of Jewish women remained lower than those of men through almost the entire nineteenth century. Even in the early twentieth century, twenty-two percent of Jewish men but only thirteen percent of Jewish women entered marriages with non-Jews. Whereas Jewish men who entered mixed marriages usually had middle-class incomes, Jewish women, by contrast, tended to marry non-Jews out of economic need or because of a lack of available male Jewish partners. And even though women’s intermarriage rates were lower than men’s, women in mixed marriages stood to lose their status in the official Jewish community, while men suffered no equivalent punishment. Male and female conversion rates similarly reflected the disproportionately high male intermarriage rates. Relatively few women converted before 1880, and when the rate increased, as it did during the years 1873–1906, women still accounted for only one quarter of all converts. In comparison with male converts, nearly double the number of women came from the lowest income categories. Rising female conversion rates appear to have coincided with the growth of secularization on the one hand, and women’s increasing participation in the workforce and ensuing encounter with antisemitism on the other. By 1912, women accounted for forty percent of all conversions. Throughout the nineteenth and early twentieth centuries, Jewish girls received an education that was consonant with social expectations for women of their class. Until the 1890s, the only form of secular education available to girls was the elementary school and non-college-preparatory secondary school. Jewish girls of all classes attended either private or public elementary schools where they learned reading, writing, arithmetic and such “feminine” subjects as art, music and literature. From mid-century on, a disproportionately high percentage of Jewish girls attended girls’ secondary schools (Höhere Töchterschule) which tended to be associated with upward mobility and higher class status. Indeed, around the turn of the century, while 3.7 percent of non-Jewish girls in Prussia attended the Höhere Töchterschule, approximately forty-two percent of Jewish girls did. Upon completing school at the age of fifteen or sixteen, middle-class girls passed their time socializing, embroidering or doing volunteer work as they waited for their families to find them a suitable husband. Even through the Imperial period, most middle-class Jewish marriages continued to be arranged either by marriage brokers or, more often, with the aid of parents and relatives. As a social institution, arranged marriage served as a means of locating Jewish marriage partners while simultaneously providing for the financial security of middle-class daughters and cementing economic alliances between families. By the end of the nineteenth century, the heavy emphasis placed on financial considerations in the search for marriage partners generated substantial criticism from within the Jewish community and particularly among young modern-minded women who wanted to choose their own life partners on the basis of romantic love. Beginning with the salon women in the eighteenth century, the decision to marry a non-Jewish man appears to have sometimes been driven at least in part by the ideal of companionate marriage. In other words, for some women, intermarriage represented not simply an act of betrayal, as it was sometimes perceived by observers, but in fact an act of independence, a rejection of a patriarchal social system that treated marriage as a financial and social transaction that was divorced from the individuals themselves. Since the nineteenth-century ideology of separate spheres consigned women to the home, those women who desired access to higher education and professional training had a particularly difficult path to navigate. For both men and women, higher education offered a means of self-improvement that facilitated German Jewish acculturation together with the possibility for personal emancipation. Yet whereas young Jewish men had been permitted to attend college preparatory high schools (gymnasia) and universities since the early nineteenth century, Jewish women had been excluded from both institutions until the end of the century. It was not until the first decade of the twentieth that German universities began admitting women. In the three years following the opening of Prussian universities to women in 1908, Jewish women already accounted for eleven percent of the female student population. By the time of the Nazi accession to power in 1933, a high proportion of Jewish women received doctorates from German universities. One of the fields of study most in demand among Jewish women, and east European Jewish women in particular, was medicine. Philosophy was also the first choice of many Jewish women since it provided the required academic preparation for a teaching certificate. As one of the few careers considered socially acceptable for middle-class women, education continued to draw Jewish women despite the antisemitic discrimination they often faced. With somewhat less frequency, Jewish women also studied the social and natural sciences and law. Despite the relative prevalence of Jewish women at universities, however, their social acceptance did not proceed apace. Like men, Jewish women encountered widespread antisemitism at the university, but their sex proved to be an added obstacle in their path toward integration. Because of the predominantly middle-class status of German Jews, fewer Jewish women were wage earners than non-Jewish women. But both single and married Jewish women did work outside the home, and they did so in growing numbers. The 1882 employment statistics for Prussia list only eleven percent of all Jewish women as part of the labor force, compared with twenty-one percent of non-Jewish women, but this figure masks the work of many more women who helped run family businesses or otherwise contributed to the household economy. In 1907, when the Prussian census included more of these invisible female workers, the employment rate was eighteen percent of Jewish women, compared with thirty percent of non-Jewish women. By the time of the Weimar Republic, with increased east European immigration, a worsening economy, and an increasing number of women working to support themselves, the gap between the Jewish and non-Jewish employment rate narrowed further, with twenty-seven percent of Jewish women now working, compared with thirty-four percent in the general population. Like Jewish men, middle-class Jewish women worked disproportionately within the commercial sector of the economy. But in contrast with native-born German women, east European immigrant working women were clustered in industrial labor, primarily in the tobacco and garment industries. In specifically low-status female occupations such as domestic service, east European immigrant women were significantly overrepresented. One of the promising new employment opportunities for Jewish and non-Jewish women at the turn of the century was social work. Formulated by women themselves as an extension of the domestic sphere, social work involved, in the words of Alice Salomon, one of the Jewish founders of modern social work in Germany, “an assumption of duties for a wider circle than are usually performed by the mother in the home” (Taylor Allen, 213–214). Jewish women seemed to flock to the profession, evident in their overrepresentation within social work training colleges. Particularly during the Weimar Republic, social work stood out as a field generally free from the mounting antisemitism increasingly being felt in other professions. Among those Jewish women who trained as social workers, some elected to work with the working class, lower middle class and east European Jewish population sectors within the Jewish community that required, in the view of their middle-class patrons, the provision of health services, job training and “moral reform.” From their roles as organizers of mutual assistance and charitable work in the eighteenth century, middle-class Jewish women became, by the Weimar period, the agents of a rationalized and “scientific” social work, one that was viewed by its practitioners as the modern-day realization of the traditional Jewish ethic of charity. As a gendered sphere of Jewish communal activity, the social arena became not only a site where those in need received assistance, but also a form of Jewish social engagement that strengthened the bonds of solidarity and cohesion among those engaged in social work. In Germany, this idea of “social motherhood” not only provided the intellectual foundation and political justification for the emergence of modern social work, but it also animated the German feminist movement from its early years until its collapse and cooptation under Hitler in 1933. Feminists’ conceptions of citizenship, rooted in distinctly organic notions of German citizenship, emphasized duties over rights and tended to define individual self-fulfillment in the context of community. Social motherhood also formed a central pillar of the German Jewish feminist movement that was founded by Bertha Pappenheim in 1904. The membership of the Jüdischer Frauenbund, which consisted primarily of middle-class married women, engaged in social work, provided career training for Jewish women, sought to combat White Slavery and fought for the equal participation of women in the Jewish community. Claiming the membership of more than twenty percent of German Jewish women, the Frauenbund became an increasingly important organization on the German Jewish scene until its dissolution by the Nazis in 1938. Middle-class Jewish women who were less interested in joining their Jewish and feminist commitments could become active in the moderate wing of the German Women’s movement, whereas working-class and east European women tended to join unions or the socialist women’s movement. Within the bourgeois women’s movement, Jewish women assumed significant leadership roles: Fanny Lewald and Jenny Hirsch gave voice to the aspirations of the movement through their writings on the “Woman Question,” while Jeanette Schwerin (1852–1899), Lina Morgenstern, Alice Salomon and Henriette Fürth became important women’s rights leaders and social workers. It has been estimated that approximately one third of the leading German women’s rights activists were of Jewish ancestry. The new democratic republic that was born amidst the catastrophe of German defeat in World War I promised Germans their first real possibility for liberal democratic governance. The constitution guaranteed equal rights to all its citizens, including full and complete equality for Jews and women. But the spirit of openness and tolerance enshrined in the constitution was quickly compromised by an eruption of virulent antisemitism that resulted in a growing economic and social exclusion of Jews, even as opportunities in some fields, such as politics and the professions, continued to expand. Weimar’s contradictory bequest to Jews—greater inclusion but also growing exclusion and intensified antisemitic rhetoric—was fueled by the ongoing economic and political instability of the period. In addition to the political instability that dogged the Republic from its inception, social and economic changes ushered in by the war also led to shifting gender roles. Many more women entered the workforce out of economic necessity and young women also sought out new professional opportunities. These and other changes in turn gave rise to the widespread perception that Germany—and German Jewry—faced an unprecedented social crisis. Rising rates of juvenile delinquency and out-of-wedlock births, the decline in the number of marriages and numbers of children born, suggested to many middle class observers that the Jewish family could neither socially nor biologically reproduce itself. Nothing embodied the social threat posed by young women to the Jewish middle-class gender norms better than the image of the sexually liberated and financially independent “New Woman,” who reputedly rejected motherhood in favor of a hedonistic urban lifestyle. What is particularly significant in the 1920s is how the identification of social crisis, as in Berlin over one hundred years before, was conceptualized largely through the lens of gender. Offering a counterpoint to the emancipated Jewish New Woman, male and female Jewish leaders placed new emphasis on the reproductive Jewish woman. Feminist leaders joined rabbis and eugenicists in calling for an increased Jewish birthrate and Jewish women’s organizations dedicated themselves to reversing Jewish women’s “self-imposed infertility” (von Ankum, 29). By reproducing Jews, women would be helping to fortify a declining Jewish community and fighting the rising tide of assimilation. In an age of assimilation, Jewish mothers had a vital role to play in the maintenance of Jewish difference itself. In the construction of a redemptive Jewish femininity that would address the challenges of assimilation, Jewish women also sought to redefine the meaning of Jewish motherhood at a time when national identity among non-Jewish Germans was growing increasingly exclusionary. According to both male and female leaders at the time, a crucial part of a Jewish mother’s task in the 1920s was to educate her children in ways that would help reduce antisemitism, while simultaneously making her family a refuge from antisemitic hostility. Shaping a new form of Jewishness that could both resist the appeal of Gentile acceptance and minimize Gentile hatred became an important aspect of Jewish “women’s work” in the 1920s. Women were thus cast both as the problem and the solution, embodying both the threat of a barren future and the promise of collective renewal. With the slide of the Weimar Republic into authoritarianism and ultimately dictatorship in the early 1930s, National Socialism signaled the end of democracy, women’s equality and Jewish emancipation in Germany. Although National Socialism targeted Jewish men and women equally, the impact of restrictive regulations, increased antisemitism and social exclusion affected Jewish men and women in ways that were often distinct. Marion Kaplan’s research on the 1930s shows how social exclusion experienced by men in the workplace appears to have had somewhat of a lesser impact than the increasing isolation from the informal social networks maintained by women. In addition, women often proved to be more attuned to the humiliations and suffering of their children. Perhaps less invested in their professional identities than their husbands, women were more willing to risk uncertainty abroad. Overall, women displayed greater adaptability than men in reorienting their expectations and their means of livelihood to accommodate new realities both at home and abroad. Ironically, it may have been women’s very subordinate status that made them more amenable to finding work that under other circumstances would have been considered beneath them. Gender roles in Jewish families also shifted as families faced new and extreme economic and social realities. Women increasingly represented or defended their husbands and other male relatives with the authorities. In addition, many more women worked outside the home than before the Nazi period and became involved in Jewish self-help organizations that had been established after Hitler’s rise to power. Some had never worked before, while others retrained for work in Germany or abroad. Although women often wanted to leave Germany before their husbands came to share their view, they actually emigrated less frequently than men. Parents sent sons away to foreign countries more frequently than daughters, and it was women, more than men, who remained behind as the sole caretakers for elderly parents. Indeed, a large proportion of the elderly population that remained in Germany was made up of women. In 1939, there were 6,674 widowed men and 28,347 widowed women in the expanded Reich. Although men and women were equally targeted for persecution and death, they were subjected to different humiliations, regulations and work requirements. Within certain types of mixed marriages, Jewish men faced greater dangers than women. In the case of childless intermarriages consisting of a Jewish woman and an “Aryan” man, the female Jewish partner was not subjected to the same anti-Jewish laws as the rest of the Jewish population. But a Jewish man with a female “Aryan” wife in such a marriage received no special privileges. With the onset of the war, German Jewish women began to suffer the kind of physical brutality that many of their husbands, fathers and brothers had endured during the 1930s. Overall, however, Jewish men were probably more vulnerable to physical attack than women. Although Jewish women who went into hiding could move about more freely and were in less danger of being discovered than men, it is speculated that fewer women than men actually went into hiding. Despite their equal status as subhuman in the eyes of the Nazis, Jewish men and women frequently labored to survive under different constraints. As was the case in other countries outside of Germany, Jewish women appear to have suffered the ultimate fate of death in disproportionately greater numbers. Even for an historical event as defining as the Holocaust, gender analysis proves a valuable means for elucidating different reactions to persecution by men and women, as well as highlighting gender-distinctive experiences of emigration, hiding and surviving in the camps. To view German Jewish history from the Enlightenment through the Holocaust from a gender perspective deepens our understanding of history in general and provides us with a richer, more complex and more inclusive picture of the Jewish past. Allen, Ann Taylor. Feminism and Motherhood in Germany 1890–1914. New Brunswick, New Jersey: 1991, 213–214; Ankum, Katharina von. “Between Maternity and Modernity: Jewish Femininity and the German-Jewish ‘Symbiosis.’” Shofar 17/4 (Summer 1999): 20–33; Baader, Maria Benjamin. “When Judaism Turned Bourgeois: Gender in Jewish Associational Life and in the Synagogue, 1750–1850.” Leo Baeck Institute Yearbook 46 (2001): 113–123; Fassmann, Irmgard Maya. Jüdinnen in der deutschen Frauenbewegung 1865–1919. New York: 1996; Freidenreich, Harriet. Female, Jewish and Educated: The Lives of Central European University Women. Bloomington: 2002; Hertz, Deborah. High Society in Old Regime Berlin. New Haven: 1988; Hyman, Paula. Gender and Assimilation in Modern Jewish History: the Role and Representation of Women. Seattle: 1992; Kaplan, Marion. Between Dignity and Despair: Jewish Life in Nazi Germany. New York: 1998; Idem. The Jewish Feminist Movement in Germany: The Campaigns of the Jüdischer Frauenbund, 1904–1938. Westport, CT: 1979; Idem. The Making of the Jewish Middle Class: Women, Family, and Identity in Imperial Germany. New York: 1991; Kaplan, Marion, ed. Geschichte des jüdischen Alltags in Deutschland. Vom 17. Jahrhundert bis 1945. Munich: 2003; Lowenstein, Steve. Berlin Jewish Community: Enlightenment, Family, Crisis 1770–1830. New York: 1994; Maurer, Trude. Die Entwicklung der jüdische Minderheit in Deutschland (1780–1933). Tübingen: 1992; Meyer, Michael, and Michael Brenner. German-Jewish History in Modern Times. New York: 1997, Vols 1–4; Quack, Sybille. Zuflucht Amerika. Zur Sozialgeschichte der Emigration deutsch-jüdischer Frauen in die USA 1933–1945. Bonn: 1995; Rahden, Till van. “Intermarriages, the ‘New Woman’ and the Situational Ethnicity of Breslau Jews from the 1870s to the 1920s.” Leo Baeck Institute Yearbook 46 (2001);125–150; Richarz, Monika. “Jewish Social Mobility in Germany during the Time of Emancipation (1790–1871).” Leo Baeck Institute Yearbook 20 (1975): 69–77; Springorum, Stefanie Schüler. “Deutsch-Jüdische Geschichte als Geschlechtergeschichte.” Transversal: Zeitschrift des David-Herzog-Centrums für jüdische Studien 1 (2003): 3–15; Usborne, Cornelie. “The New Woman and Generational Conflict: Perceptions of Young Women’s Sexual Mores in the Weimar Republic.” In Generations in Conflict: Youth Revolt and Generation Formation in Germany, 1779–1968, edited by Mark Roseman, 137–163. New York: 1995; Volkov, Shulamit. Die Juden in Deutschland 1780–1918. Munich: 1994; Idem. “Jüdische Assimilation und Eigenart im Kaiserreich.” In Jüdisches Leben und Antisemitismus im 19. und 20. Jahrhundert, edited by Shulamit Volkov. Munich: 1990, 131–145; Werthheimer, Jack. Unwelcome Strangers. New York: 1987; Zimmermann, Moshe. Die deutschen Juden, 1918–1945. Munich: 1997. How to cite this page Gillerman, Sharon. "Germany: 1750-1945." Jewish Women: A Comprehensive Historical Encyclopedia. 1 March 2009. Jewish Women's Archive. (Viewed on October 22, 2016) <http://jwa.org/encyclopedia/article/germany-1750-1945>.
Astronomers using NASA’s Hubble Space Telescope have discovered what they believe to be some of the greatest evidence yet for the presence of a rare type of “intermediate-sized” black hole lurking in the heart of the nearest globular star cluster to Earth, about 6,000 light-years distant. What are the formation, distribution, and rarity of intermediate-mass black holes? Almost all black holes appear to come in two sizes, similar to strong gravitational pits in the fabric of space: small and gargantuan. Our galaxy is thought to be strewn with 100 million tiny black holes (many times the mass of our Sun) formed by exploding stars. The cosmos is teeming with supermassive black holes, which are situated in the centers of galaxies and weigh millions or billions of times the mass of our Sun. An intermediate-mass black hole, weighing between 100 and 100,000 solar masses, is a long-sought missing link. How would they form, where would they congregate, and why do they appear to be so uncommon? Using a variety of observational approaches, astronomers have detected more probable intermediate-sized black holes. Three of the finest possibilities — 3XMM J215022.4055108, discovered by Hubble in 2020, and HLX-1, discovered in 2009 — live in dense star clusters on the edges of neighboring galaxies. Each of these hypothetical black holes has tens of thousands of suns in mass and may have formerly resided in the centers of dwarf galaxies. NASA’s Chandra X-ray Observatory has also aided in the discovery of numerous probable intermediate black holes, including a large sample in 2018. Much closer to home, a number of probable intermediate-sized black holes have been discovered in dense globular star clusters around our Milky Way galaxy. For example, Hubble researchers reported the possible presence of an intermediate-mass black hole in the globular cluster Omega Centauri in 2008. These and other intermediate-mass black hole discoveries remain inconclusive and do not rule out alternate hypotheses for a variety of reasons, including the need for further data. Hubble’s unique capabilities have now been utilized to hone in on the core of the globular star cluster Messier 4 (M4), allowing for more precise black-hole hunting than prior efforts. “You can’t do this kind of science without Hubble,” said Eduardo Vitral, lead author of an article to be published in the Monthly Notices of the Royal Astronomical Society. What Vitral’s team discovered and what was its significance? Vitral’s team discovered a probable 800 solar-mass intermediate-sized black hole. Although the alleged object cannot be seen, its mass can be determined by observing the motion of stars caught in its gravitational field, similar to bees swarming around a hive. Measuring their movement needs time and precision. This is where Hubble achieves something that no other modern telescope can. Astronomers examined 12 years of Hubble M4 data and resolved pinpoint stars. His team believes the black hole in M4 could be 800 times the mass of our Sun. Alternative possibilities for this object, such as a compact center cluster of unresolved stellar remains like neutron stars or smaller black holes revolving around one other, are ruled out by Hubble’s observations. “We are confident that we have a very small region with a large concentration of mass.” “It’s about three times smaller than the densest dark mass we’ve found in other globular clusters,” Vitral added. “When we consider a collection of black holes, neutron stars, and white dwarfs segregated at the cluster’s center, the region is more compact than what we can reproduce with numerical simulations.” They are not capable of forming such a dense concentration of mass.” Credits: NASA’s Goddard Space Flight Center; Lead Producer: Paul Morris What is the nature of the central mass in the globular cluster and its impact on stellar motions? A collection of closely packed objects would be dynamically unstable. If the object isn’t a single intermediate-sized black hole, the observed stellar motions would require an estimated 40 smaller black holes squeezed into an area only one-tenth of a light-year across. As a result, they would merge and/or be expelled in an interplanetary pinball game. “We measure the motions and positions of stars and apply physical models to try to reproduce these motions.” “We end up with a measurement of a dark mass extension in the center of the cluster,” Vitral explained. “The stars move more randomly as they get closer to the central mass.” And, as the center mass increases, so do the stellar velocities.” Because intermediate-mass black holes in globular clusters have been so difficult to find, Vitral warns, “While we cannot completely confirm that it is a central point of gravity, we can show that it is very small.” It’s too little for us to explain anything other than a solitary black hole. Alternatively, there could be a stellar mechanism that we are simply unaware of, at least in terms of present physics.”
Weather forecasting is the classic inexact science, relying on the complex mutual interactions of wind, currents, precipitation, tides, humidity, and temperature variations, and a million other variables across a planet that's rotating on its axis, revolving around its heat source, and tilted with regard to its plane of revolution. To say forecasting the weather is tricky is putting it mildly indeed. In fact, it was while working on weather prediction that mathematician Edward Lorenz began to conceive Chaos Theory, the mathematical theory which says some systems, highly sensitive to initial conditions, are simply too complex to be predictable over the long term. Weather is the poster child for chaos theory: Its most famous idea, the "Butterfly Effect," posits that a butterfly flapping its wings in South America could have far-ranging repercussions for future weather around the world, like a tornado in Texas. But weather forecasting is too important to be left to chance, and humans have been doing our best to predict the weather throughout history. How much closer have our sophisticated technology and global communication advances brought us to useful predictability? We decided to investigate. Weather prediction stayed amateurish and local until the telegraph became widely available in the middle of the 19th century, letting local weather forecasters share their observations from all over the world, displaying this parallel information in "synoptic" weather charts, so people could see the weather occurring at the same time in different places. For the first time, predictions of large-scale moving fronts was fact-based and reliable. The first U.S. national weather forecasting service was formed in 1870, in response to the need for storm warnings on the Great Lakes; by the 1930s weather balloons were dotting the sky, using radio waves to relay high-altitude readings and general weather patterns back to forecasters on the ground. Today, a network of more than a thousand operational weather satellites feeds the forecasters of the world, carrying a wide array of sensors utilized in different ways. Geostationary satellites circle the Earth from west to east, in the equatorial plane at an altitude of about 22,500 miles, matching the pace of the Earth's rotation so as to stay above the same spot on the equator. These satellites provide a continuous view of the hemisphere from their position, feeding invaluable realtime information about the rapid condition shifts that occur during thunderstorms, hurricanes, and fronts. Polar Orbiters are low-flying, circling the Earth in roughly north–south orbits at much lower altitudes (from 310 to 620 miles above the surface). Each orbiter flies the same path twice a day, and their relatively low altitude provides extremely high-resolution images and extremely specific atmospheric conditions, making predictive patterns far more accurate. Radar began to emerge as a meteorological tool in the 1930s, providing excellent views of raindrops at certain wavelengths. At the local level, radar made it possible to study the evolution of thunderstorms, "see" their precipitation structure, and track their progress. Today, vitually all tornadoes and severe thunderstorms in the United States and around the world are monitored by radar, charting their characteristics and severity. Today's high-end weather prediction systems typically employ "Doppler Radar," which takes advantage of the Doppler effect to track storms and other moving phenomena. As a Doppler radar antenna turns, it emits short bursts of radio waves called "pulses" through the atmosphere. When the wave hits moving objects, it reflects back toward the source, which also contains a receiver as well as the original transmitter; by measuring the shift in phase between a transmitted pulse and the received echo it creates, the moving object's radial velocity (the movement of the target directly toward or away from the radar) can be calculated. In the aggregate, this data can be expressed as those familiar swirling depictions of weather patterns, and allows detailed analysis of their growth and movement. PROCESSING THE DATA From the '60s through today, the increasing power and ubiquity of high end computer systems has allowed a quickly improving ability to process these vast data sets. Numerical weather prediction (NWP) involves constructing a three-dimensional grid of the atmosphere using data on the atmosphere's current state, and using a mathematical model to predict its evolution into future states. Taking readings on the temperature, humidity, barometric pressure, wind speed and direction, and precipitation for each point on that grid, and repeating that process every 10 minutes or so, provides a stream of data that feeds into mathematical algorithms that seek to predict future states. These algorithms' success is later analyzed, and revised to improve predictive quality over time. Those methods have become fairly reliable in short-term forecasting of one to three days: If the weatherman calls on Thursday for weekend rain, cancel that picnic. But for longer-range forecasting, chaos theory suggests that this extremely complex system, so dependent on early conditions, starts to diverge quickly, and the fifth or sixth day of that "seven-day forecast" is significantly less reliable than the third and fourth. Long-range forecasters prefer a climatological approach, focusing on a broad weather picture over a period of time rather than attempting to pinpoint daily details. Forecasters must then employ a process known as "ensemble forecasting" that takes into account multiple models, which can be inconsistent over time and therefore unpredictable in the specifics. They try to predict the departures from normal conditions over a given period of time, and display these on an "anomaly map." Of course people are people, and demand to know whether it'll rain at their wedding two weeks from Sunday, even if it's implausible. The weather world tries to comply, and that's how you get to Doppler offering a "25-day forecast." Josh Rosenberg did a study checking the validity of his region's, and as you can see from the chart below, it's not a pretty picture; with the very best technology has to offer, predictions that we would call "good" don't begin to emerge until about five days prior. So while time and technology have greatly enhanced the accuracy of forecasting, the pace of advance has slowed to a crawl of minor incremental gains. Like the economy, another extremely complex system, the weather has yielded its secrets about the here and now, but this extremely complex system is resisting our efforts to accurately predict the future. More from World Science Festival... - Mer-monkeys and guitar nipple: 10 curious scientific hoaxes - World's first solar-powered circumnagivation ready for takeoff - As time goes by (in our heads)
From: Astrobiology Magazine Posted: Monday, November 10, 2008 By Michael Schirber Heat-loving organisms live where the water is hot but the gene pool is shallow. Genetic analysis has shown that so-called thermophiles have fewer mutations in their protein-coding genes than do their microbial cousins that live at room temperature. This seems to imply that the opportunities to evolve decrease as temperature increases. To confirm this, Jonathan Silberg of Rice University and colleagues will explore how resilient proteins are to mutations at different temperatures. The results could have implications for theories that claim life arose near hydrothermal vents. It is not surprising that high temperatures are tough on living things. If exposed to extreme heat, most cells will die because their outer membrane becomes unstable or because essential molecules, like proteins and DNA, cease to function properly. Proteins are especially sensitive to temperature: they become denatured (i.e. change their shape) when over-heated. "Most proteins are marginally stable at their natural environmental temperature, which is to say that they aren't all that much more stable than they need to be," says Jesse Bloom of California Institute of Technology in Pasadena, CA. "At one time, this might have been thought to imply that these proteins just couldn't get any more stable." However, scientists have since discovered thermophilic organisms whose proteins work just fine at temperatures as high as 121 degrees Celsius (249.8 degrees Farhenheit). This begs the question: can proteins be stable at even higher temperatures? Or do thermophiles mark the temperature boundary for protein-based life? Silberg thinks his group can answer this question by studying the potential for proteins to evolve as a function of temperature. "Our basic question is: how hard is it for Nature to design a protein at high temperature?" Silberg says. As part of NASA's Astrobiology: Exobiology and Evolutionary Biology program, Silberg's group will compare the effects of mutations on two versions of an essential type of protein found in both thermophiles and normal temperature mesophiles. Proteins perform a multitude of roles in living organisms. As enzymes, they help break down food; as receptors, they receive signals from other cells; as connectors, they provide structure. A single species will have many thousands of different proteins, each characterized by a sequence of amino acids linked together in a chain. These chains fold up on themselves to form a complex three dimensional shape. Often it is this shape that allows a protein to do its job. However, the shape is not uniquely tied to one chain of amino acids: there can be several different chains that all fold up into the same shape and, therefore, all do roughly the same thing. This redundancy is thought to be crucial to the evolution of proteins because it allows for a certain amount of random mutations. Imagine an organism relies on protein X to do some vital task inside its cells. One day a mutation occurs in the gene that codes for protein X, resulting in a slightly different amino acid sequence. Two things can happen. One, the mutated protein X folds in a different way and is therefore useless. The organism ends up dying or is unfit to compete with its neighbors. The other possibility is that the mutated protein X retains the original shape and therefore continues to perform the same task as before -- a "shape-preserving mutation." Each protein has a whole network of shape-preserving mutations associated with it. "Over time, organisms drift across these networks in order to find some slight improvement," Silberg says. Those genetic mutations that confer some advantage will spread throughout the population. "That's how adaptive evolution works," Silberg says. In earlier work, Silberg and colleagues devised a model that predicts how many shape-preserving mutations are associated to a particular protein. The model is based on thermodynamic parameters that give the probability that a certain amino acid sequence has a particular shape. "My idea comes from a physics perspective," Silberg says. "Proteins are polymer chains in which each bead can be one of 20 different amino acids, so what happens if you change one of the beads?" Although most of the work is done with computer simulations, the researchers have tested the accuracy of their model by synthesizing proteins in the lab with various amino acid substitutions. The results provide a picture of the evolutionary playing field on which protein-coding genes can explore possible improvements. Gene pool put on simmer In this new project, Silberg's group plans to look for the first time at the number of shape-preserving mutations available to thermophiles. Their hypothesis is that the evolutionary playing field shrinks - that there are fewer possibilities for viable mutations --as the temperature goes up. Previous work has shown that the protein-coding genes of thermophiles have an unusually low number of a particular type of mutations. The cause of this may not be due to temperature, because reduced genetic diversity can arise in a number of ways. Silberg and his collaborators want to directly examine a protein at high temperature to see if it is less tolerant to slight changes in its amino acid configuration. This experiment could confirm whether or not genetic diversity is correlated with temperature. The protein the group has chosen to study is adenylate kinase (AK), which regulates the energy inside a cell and is essential to all living things. The researchers will look at AK in two bacteria: a thermophile (Thermus thermophilus) and a mesophile (E. coli). The AK protein is similar in both organisms, but there are distinct differences due to the fact that they operate at different temperatures (80 and 40 degrees C, respectively, (176 and 104 F)). The researchers will determine-through a combination of lab work and computer simulations-which amino acid substitutions retain the original shape of each protein at their respective temperatures. By comparing the number of shape-preserving substitutions, the team will be able to say how temperature affects protein evolution. "The research proposed by Dr. Silberg is potentially interesting because it will indicate whether the proteins of thermophiles are anywhere near the maximum thermostability that is evolutionary achievable," says Bloom, who is not involved with this study. If thermophilic proteins are just barely within the evolutionary limits, this might argue against hydrothermal vents and/or hot volcanic pools being home to the first life on Earth (a popular theory due to the fact that thermophiles are closely related to life's earliest common ancestor). If high temperatures prove to be less forgiving of the genetic mutations that are necessary for evolution, then it might make more sense that life started in cooler environs and only later migrated into hot places. Furthermore, if Silberg and company do find a temperature dependence to protein mutation tolerance, the results might be extrapolated to higher temperatures to locate at what point the evolutionary playing field shrinks to zero."This will tell us where proteins no longer make for a good polymer," Silberg says. If life were to exist on an exoplanet hotter than this maximum protein temperature, those extra-terrestrials would likely need a different "do-it-all" molecule to replace proteins. // end //
The Typographic Scale – Harmony in FontsDate: 6/1/2020, Author: G.P. The classic typographic scale is a method that refers to the way font sizes are progressing. Typography, from print media to the web, has always needed a way to present text in a consistent manner. Text should always be presentable and easy to read. So, that’s where typographic scales come into play. They offer us a system that allows us to use text harmonically. They actually resemble the way musical notes and musical scales are working. Here is the classic typographic scale that can be found in many, commonly used, applications: 6 7 8 9 10 12 14 16 18 21 24 28 32 36 42 48 55 63 73 84 96 You will notice that some numbers are in bold. You can also easily notice that every one of those bold numbers is a multiple of the previous number and the number 2. So, the number 2 is the ratio of our scale. As you can see, there are always exactly 5 numbers from the first bold number to the next one. This is the interval of the scale. So, lets summarize: A typographic scale consists of 2 basic properties: - A ratio (r). Every size must be multiplied by the ratio to find the next harmonic size. For example, the classic typographic scale has a ratio of 2. A ratio of 1.5 (which is called a perfect fifth in music), means that any font size is exactly 1.5 times bigger than the one before it. We call it a perfect fifth because in a musical scale with a ratio 1.5, there are always 5 notes in the sequence. In the same way we can also use the Golden Ratio (1.618) to create harmonic typographic scales. - An interval (s). The number of sizes from the first element to the last element of a sequence. For example, if we have this scale: 12, ..., …, …, …, 24, ..., …, …, …, 48 ..., …, …, …, we have a ratio r=2 and a number of sizes s=5. So, if i is the ith element of the scale and n is the size number from which we want to know the ith harmonic size, we can use the basic formula for the frequencies of the notes of the equal tempered scale (a scale with the same ratio): xi = nri/s Thus, if we want to know the 2nd size from 12 we can find it using the previous formula: x2 = 1222/5 = 1220,4 = 121,319 = 15,82 ≈ 16 Similarly, the 3rd size from 12 is: x3 = 1223/5 = 1220,6 = 12*1,515 = 18,18 ≈ 18 Our scale now is: 12, 16, 18, …, …, 24, ..., …, …, …, 48 ..., …, …, …, We can continue to fill all the elements in the same way we found the 2nd and the 3rd size. This will give as a harmonic size scale. How to use typographic scales on the web A simple way to use a typographic scale on the web is to find the basic size of the body. You can use either pixels or ems (there is a difference that we will analyze in a next article about Responsive Typography). After that you can assign the next harmonic size to the smaller header (h6, h5, h4, h3, h2, h1). You could also revert the process and start with the largest size. Finally, some ratios have been proven more harmonious than others. Here is a list of some scales you could experiment with: You can find the full scales here - https://type-scale.com/ Of course, you can always experiment using your own values. The Typographic Scale by Spencer Mortensen Equal Temperament - https://en.wikipedia.org/wiki/Equal_temperament The Golden Ratio - https://en.wikipedia.org/wiki/Golden_ratio Perfect Fifth - https://en.wikipedia.org/wiki/Perfect_fifth
Search Within Results Common Core: Standard Common Core: ELA Common Core: Math - In Topic C, students graph logarithmic functions, identifying key features (F-IF.4, F-IF.7) and discover how the logarithmic properties are evidenced in the graphs of corresponding logarithmic... - Topic D opens with a hands-on simulation and modeling activity in which students gather data and apply the analysis of Lesson 22 in Topic C to model it with an exponential function (A-CED.2, F-LE.5... - In Topic A, students both review what they already know about exponential expressions and functions with integer exponents, and extend the meaning of an exponential expression to allow for first... - In Topic A, students explore arithmetic and geometric sequences as an introduction to the formal notation of functions (F-IF.A.1, F-IF.A.2). They interpret arithmetic sequences as linear functions... - In Topic D, students apply and reinforce the concepts of the module as they examine and compare exponential, piecewise, and step functions in a real-world context (F-IF.C.9). They create equations... - Topic A focuses on the skills inherent in the modeling process: representing graphs, data sets, or verbal descriptions using explicit expressions (F-BF.A.1a) when presented in graphic form in Lesson... - Tables, graphs, and equations all represent models. We use terms such as “symbolic” or “analytic” to refer specifically to the equation form of a function model; “descriptive model” refers to a...
Describing Data There are two ways to describe a set of data: • Graphically Dot plot – Bar Graph – Histogram – Boxplot • Numerically-using a single number to describe the relationship of the data • Measures of Center: Mean and Median review these descriptions in your notes • Measures of Spread- Today we will focus on the numerical descriptions Measures of Spread describe how much values typically vary from the center These measures are: • Range – a description of how far apart the distance from the highest to lowest data pieces; found using highest minus lowest data • Interquartile Range (IQR) – a description of the middle 50% of the data ; found using Q3 – Q1 • Mean Absolute Deviation briefly view description in your notes and then continue with the PPT. Now let’s explore “deviations from the mean” as a way to determine how accurately the mean may describe “typical”. Thinking about the Situation Consider the following test scores: Who is the best student? How do you know? Take a few minutes to decide with your partner So with all students having the same mean, let’s see if we can dig a little deeper in our comparison. The MEAN ABSOLUTE DEVIAITON (MAD) will help us. Let’s take a few minutes to explore MAD. Mean Absolute Deviation (MAD) Add to notes STEP 1: Find the mean STEP 2: Subtract the mean from each piece of data STEP 3: Find the absolute value of each difference STEP 4: Find the mean of the new differences (deviations) NOW WE WILL TRY THIS METHOD. Mean Absolute Deviation Write the list of numbers shown on the number line and then find the Mean Absolute Deviation 0 1 2 3 4 5 6 7 8 9 10 11 STOP AND COMPLETE CHART on NOTES So what exactly is deviation? -4 -3 +5 -1 +3 0 1 2 3 4 5 6 7 8 9 10 11 (-4) + (-3) + (-1) = -8 (+5 ) + (+3) = +8 Mean Absolute Deviation -3.2 +3.2 0 1 2 3 4 5 6 7 8 9 10 11 Notice that our Mean Absolute Deviation or MAD was 3.2 and most of our original data does fall within plus or minus 3.2 points of the mean of 5. Highlight the following in your notes: A low mean absolute deviation indicates that the data points tend to be very close to the mean and not spread out very far so the mean is an accurate description of “typical”, and a high mean absolute deviation indicates that the data points are spread out over a large range of values. Now take a few minutes to go back to our original question about the best student. Find the MAD score for each student and then make a decision based on all of your data about the best student. Be prepared to discuss. Now you try the CLASSWORK found on your handout. When you and your partner have completed the work individually, check to see if you agree. We will review this section together in a few minutes
Reasoning – Reasoning refers to the process by which people generate and evaluate arguments and beliefs. Philosophers have long distinguished between two kinds of reasoning: inductive and deductive. We examine each separately here, and will then explore one of the most powerful mechanisms people use to make inferences, particularly about novel situations: reasoning by analogy. Inductive reasoning is reasoning from specific observations to more general propositions. An inductive conclusion is not necessarily true because its underlying premises are only probable, not certain. For example, say you asked a friend who appeared to be quite upset if they were feeling ok, and they replied ‘yes I’m fine’, inductive reasoning could lead you to the conclusion ‘if my friend says she is fine, then she must be fine’. It is simply reasoning made by observation. Another example of inductive reasoning would be a child who believes that ‘if Santa can climb down the chimney, so can the boogey man!’ Nevertheless, inductive reasoning forms a large chunk of our day to day reasoning. If someone raises their voice we reason that they must be angry, if someone looks sad, we reason that something must be wrong, if everything appears the same as it did yesterday, we reason that everything is the same as yesterday. Deductive reasoning is logical reasoning that draws a conclusion from a set of assumptions, or premises. In contrast to inductive reasoning, it starts off with an idea rather than an observation. In some ways, deduction is the flipside of induction: whereas induction starts with specifics and draws general conclusions, deduction starts with general principles and makes inferences about specific instances. For example, if you understand the general premise that all dogs have fur and you know that your next door neighbour just bought a dog, then you can deduce that your neighbour’s dog has fur, even though you haven’t seen it yet. This kind of deductive argument is referred to as a syllogism. A syllogism consists of two premises that lead to a logical conclusion. If it is true that: A) All dogs have fur and B) The neighbour’s new pet is a dog. Then there is no choice but to accept the conclusion that: C) The neighbour’s new dog has fur. Unlike inductive reasoning, deductive reasoning can lead to certain rather than simply probable conclusions, as long as the premises are correct and the reasoning is logical. Deductive reasoning seems as though it would follow similar principles everywhere. However, recent research suggests that Eastern and Western cultures may follow somewhat different rules of logic – or at least have different levels of tolerance for certain kinds of inconsistency. The tradition of logic in the West, extending from Ancient Greece to the present, places an enormous emphasis on the law of non-contradiction: two statements that contradict each other cannot both be true. This rule is central to solving syllogisms. Contrastingly however, in the East, people view contradictions with much more acceptance, and often believe them to contain great wisdom. Take for example Zen Koans (statements or questions) that offer no rational solutions, and are often paradoxical in nature, such as: if a tree falls in the forest, and no one is around to hear it, does the tree make a sound?’ Western thought aims to try and resolve contradiction by using logic, while people in the East focus instead on the truth that each statement provides – relishing, rather than resolving paradox. A humorous example of this difference can be found in an episode of South Park, where Stan and Kyle are trying to locate and destroy ‘the heart of Wal-Mart’ so it will lose it’s power over South Park’s inhabitants. In this exchange of dialogue, Stan and Kyle represent Western rational thought, while the embodiement of Wal-Mart represents Eastern paradoxical thought: STAN: We don’t want your store in our town; we’ve come to destroy you! KYLE: Where’s the heart? WAL-MART: To find the heart of Wall-Mart, one must first ask oneself, “Who is it that asked the question?” STAN: Me. I’m asking the question. WAL-MART: Ah, yes, but who are you? STAN: Stan Marsh. Now where is the heart? WAL-MART: Ah, you know the answer, but not the question! KYLE: The question is: ‘Where is the heart?’ In this exchange of dialogue you can see that Stan/Kyle, the Western rational thinkers, have a low tolerance for all of Wal-Mart’s philosophical questions and are answering them with straight forward logic, even though these are the not the answers that Wal-Mart (Eastern thought) seek. To Eastern thought, the question is generally more important than the answer, just as the journey is viewed as being more important than the destination. Reasoning by Analogy – Analogy is one of the single most powerful reasoning devices we have, and we use it a lot in everyday conversation to explain new situations. Analogical reasoning is the process by which people understand a novel situation in terms of a familiar one. For example, someone who has never done heroin might ask someone who has what the sensation is like; considering the person asking the question has never tried it, the heroin user would have to use analogy to create a comparison. He might say, it is like sinking into a hot bath on a freezing cold day. Since the person who asked the question can relate to the analogy (let’s assume they both live in a very cold country), through analogical reasoning he can deduce what it must feel like to be on heroin. Similarly, someone with schizophrenia could use analogical reasoning to suggest that having schizophrenia is like living in a nightmare, and not being able to wake up. A key aspect of analogies of this sort is that the familiar situation and the novel situation must each contain a system of elements that can be mapped onto one another. For an analogy to take hold, the two situations need not literally resemble each other, however, the elements of the two situations must relate to one another in a way that explains how the elements of the novel situation are similar to that of the unfamiliar situation.

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